Journal of Mathematical Psychology 44, 408463 (2000) doi:10.1006jmps.1999.1260, available online at http:www.idealibrary.com on Stochastic Dynamic Models of Response Time and Accuracy: A Foundational Primer Philip L. Smith University of Melbourne, Parkville, Victoria, Australia A large class of statistical decision models for performance in simple information processing tasks can be described by linear, first-order, stochastic differential equations (SDEs), whose solutions are diffusion processes. In such models, the first passage time for the diffusion process through a response criterion determines the time at which an observer makes a decision about the identity of a stimulus. Because the assumptions of many cognitive models lead to SDEs that are time inhomogeneous, classical methods for solving such first passage time problems are usually inapplicable. In contrast, recent integral equation methods often yield solutions to both the one-sided and the two-sided first passage time problems, even in the presence of time inhomogeneity. These methods, which are of particular relevance to the cognitive modeler, are described in detail, together with illustrative applications. 2000 Academic Press Theories of how human subjects make decisions in simple perceptual and cognitive tasks often propose some form of sequential sampling mechanism to explain the patterns of response time (RT) and accuracy that are observed in empirical data. Indeed, for some researchers, the study of such mechanisms addresses one of the most fundamental questions in psychology, namely, how the central nervous system translates perception into action and how this translation depends on the intentions and expectations of the individual. From this perspective, the study of simple decisions helps illuminate the nexus of perception, thought, and action and thus has implications for our understanding of the diverse perceptual and cognitive phenomena in which such decisions are involved. Like signal detection theory (SDT) (Green 6 Swets, 1966), the theory of sequential sampling mechanisms starts from the premise that simple perceptual and cognitive The idea for this article emerged from a seminar on response time models at Indiana University during the fall of 1996 in which I participated, and I thank Richard Shiffrin, Jerry Busemeyer, and Jim Townsend for encouraging me to write it. Background work was carried out during a sabbatical at the Institute for Mathematical Behavioral Sciences at the University of California, Irvine, in the same year, and I thank Duncan Luce for the support and hospitality of the Institute during this time. My thanks also to Jerry Busemeyer and Michael Rudd for their careful reading of an earlier version of the manuscript. Preparation of this article was supported in part by Australian Research Council Grant A79802778. Address correspondence and reprint requests to Philip L. Smith, Department of Psychology, University of Melbourne, Parkville, Vic. 3052, Australia. E-mail: philipsherman.psych.unimelb.edu.au. 0022-249600 35.00 Copyright 2000 by Academic Press All rights of reproduction in any form reserved. 408 STOCHASTIC DYNAMIC MODELS 409 decisions are statistical in nature. That this is so follows from the widely held assumption that sensory and cognitive systems are inherently noisy. Thus, a formal model of simple decisions typically consists of a set of representational assumptions, which specify how stimulus properties are encoded statistically in the central nervous system, and a set of process assumptions, which specify how this noisy information is used to arrive at a decision. Also in common with SDT, the idea that sensory representations are noisy and time-varying leads naturally to the view that simple decisions involve a smoothing or filtering operation, and this in turn leads to the presumption that they are performed by an averaging or integration device. Where SDT and the theory of sequential sampling models diverge is in their assumptions about the accrual of stimulus information. Whereas SDT assumes a fixed sampling interval, sequential sampling models assume that the interval is variable and depends on the statistical properties of the signal itself. Rather than taking a sample of stimulus information of predetermined size, such models assume that the decision mechanism samples until a criterion quantity of information needed for a response is obtained. Because this quantity varies with the statistics of the sample, the time needed to acquire it is also variable. Typically, this acquisition time is identified with the decision time component of RT. Among the many researchers who have investigated the formal properties of these models are Ashby (1983), Audley and Pike (1965), Busemeyer and Townsend (1992, 1993), Diederich (1995, 1997), Edwards (1965), Emerson (1970), La Berge (1962), Laming (1968), Link (1975, 1978, 1992), Link and Heath (1975), Luce and Green (1972), Pike (1966, 1968), Ratcliff (1978, 1981), Smith and Vickers (1988, 1989), Townsend and Ashby (1983), Vickers (1970; Vickers, Caudrey, 6 Willson, 1971), and Viviani (1979a, 1979b). Reviews of this literature may be found in Vickers (1979), Townsend and Ashby (1983), and Luce (1986). In practice, the study of a given sequential sampling model reduces to the study of a stochastic process, or processes, that represents the accumulated information available to the decision mechanism at a given time. This process (for the moment defined only in the singular) will be denoted X(t). Formally, X(t) is a random variable defined on the probability space of all possible sequences of accumulated information at time tthe term ``information'' being used here in a sense that is only loosely related to its more technical, statistical meaning. Specifically, at each time t, t # T, where T denotes the set of indices for which the process is defined, X(t) takes on a random value x, x # X. The set X, of possible values of accumulated information, is the state space of the process. The set of possible time indices, T, and the set of possible information states, X, may independently be either discrete or continuous, depending on the assumptions of the model. This article will be concerned exclusively with continuous-time, continuous state-space models. The interested reader is referred to the preceding references for other possibilities. Two steps are required to specify the quantitative properties of a model, once the set of process assumptions that determine X(t) are given. The first step is to specify the probabilistic character of the process X(t). These properties are expressed via its transition distribution, F(x, t | y, {), which is defined as F(x, t | y, {)=P[X(t)x | X({)= y]. (1) 410 PHILIP L. SMITH This distribution gives the probability that the accumulated information at time t is less than or equal to x, given that its value at some earlier time { was y. The second step is to solve the first passage time problem for X(t). Usually, the statistics of X(t) most of interest in applications are the probability that the accumulated information will eventually reach or exceed a particular level, together with the time required for this to happen. Depending on the application, the one-sided or twosided first passage time problem may be of greatest interest. The one-sided problem involves a single random variable T, defined as T=inf [t: X(t)a]; X(t 0 )<a. (2) That is, T is the time at which the accumulated information first reaches the value awhere it is assumed that the process starts at time t=t 0 with an initial value X(t 0 )=x 0 <a. We will take t 0 =0 and x 0 =0 without further comment, except where greater clarity is obtained by indicating the values of the initial conditions in an explicit way. Depending on the assumptions of the model, x 0 may either be fixed or random with a prescribed distribution. For x 0 fixed we stipulate that P[X(t 0 )=x 0 ]=1. Mathematically, a is an absorbing barrier or boundary for X(t), which is identified psychologically with the decision criterion adopted by the subject in an experimental task. The two-sided problem involves a pair of random variables, T 1 and T 2 , defined as T 1 =inf [t: X(t)a 1 ; X({)a 2 for all {<t] T 2 =inf [t: X(t)a 2 ; X({)a 1 for all {<t], (3) a 2 <X(t 0 )<a 1 . The variables T 1 and T 2 denote, respectively, the time required for the accumulated information first to exceed a 1 or to fall below a 2 , given that the other boundary has not been crossed already. Again, the values a 1 and a 2 may be interpreted as decision criteria associated with competing responses. To ensure that T 1 and T 2 are well defined, we adopt the convention that the infimum of the empty set is infinity. With this convention, the definition (3) implies that either T 1 = or T 2 = for each realisation of X(t). The statistics of interest for the one-sided problem are P[T<], the probability that X(t) exceeds a in finite time, and G T (t), the associated first passage time distribution. If T is not finite with probability one, then G T (t) will be a defective distribution; that is, its total probability mass will be less than one. The corresponding statistics for the two-sided problem are P[T 1 <T 2 ], the probability that the process crosses the boundary a 1 before crossing the boundary a 2 , and G 1(t) and G 2(t), the first passage time distributions for the random variables T 1 and T 2 . The functions G 1(t) and G 2(t) represent the joint distributions of the events that a boundary crossing occurs at or before time t and that the first boundary crossed is a 1 or a 2 , respectively. If the random variable T is defined as T=min(T 1 , T 2 ), the time of the first boundary crossing, independent of the boundary involved, then G T (t)=G 1(t)+G 2(t), STOCHASTIC DYNAMIC MODELS 411 because first boundary crossings at a 1 and a 2 are mutually exclusive events. For the continuous-time processes of interest here, T 1 , T 2 , and T possess density functions, denoted g 1(t), g 2(t), and g T (t), respectively. Specifically, g T (t)=(ddt) G T (t), with the other densities being defined similarly. When we wish to show the dependency of these densities on the absorbing boundaries and on the initial conditions we will write g 1(a 1 , t | x 0 , t 0 ) and g 2(a 2 , t | x 0 , t 0 ) for the two-barrier case and g(a, t | x 0 , t 0 ) for the single-barrier case. To provide some intuitive content for these probabilistic ideas, Figs. 1 and 2 show some simple model schemes, in which the preceding quantities arise naturally from theoretical considerations. The model in Fig. 1 provides a possible framework for modeling simple RT, GoNo Go RT, or threshold increment detection. 1 Conceptually, it consists of two main parts: an encoding stage and a decision stage. The encoding stage, shown in the box on the right, represents the properties of some particular sensory system, memory system, or systems of both kinds acting in concert. Its output is a time-dependent information function, denoted +(t). This function represents the instantaneous, encoded value of those attributes of the stimulus that are relevant to the decision task at hand. Although the exact substrate of this function is usually left unspecified, it is typically assumed to correspond to the level of activation in some neural pathway or pathways. To account for the behavioral variability that is ubiquitous in simple decision tasks, the information function is subject to continuous statistical perturbation. This perturbation is represented in the figure by the function W(t). The second part of the model, the decision stage, is shown in the box on the right side of Fig. 1. To make a decision about the stimulus, the instantaneous values of the noise-perturbed information function are summed or accumulated until the criterion, a, is exceeded. The accumulated information is represented by the stochastic process X(t). As described previously, the decision time in the model is T, the first passage time of X(t) through the level a. Versions of this model have been considered by, among others, Diederich (1995), Emerson (1970), Pacut (1977, 1980), and Smith (1995). To model choice RT, or discrimination between similar stimuli, the scheme shown in Fig. 1 can be extended in various ways. Two possibilities are shown in Fig. 2. In both, the information function is assumed to take on both positive and negative values. In cognitive models, this is often accomplished by the device of comparing sampled values of the information function to a decision referent c, prior 1 Threshold increment detection can be modeled using either a single criterion or a pair of criteria. Dual-criterion models (e.g., Link, 1978) associate a criterion with each of the responses ``Signal'' and ``Noise'' and assume that information is sampled until one of the two criteria is exceeded. Single-criterion models associate a criterion with the signal response only and assume that a noise response is emitted by default if the signal criterion is not exceeded in some predetermined period of time. Well-behaved psychometric functions can be obtained with single criterion models only if the asymptotic probability of a signal response is less than unity; that is, if P[T<]<1. Typically, this means that single-criterion models will be appropriate only when stimulus duration is limited, as occurs when stimuli are backwardly masked. When weak, response-terminated stimuli are used, dual-criterion models are needed to obtain psychometric functions with the required properties. These properties can often be characterized by considering the boundary behavior of the accumulation process (e.g., Karlin 6 Taylor, 1981, pp. 226250) to ascertain whether absorption at a boundary in finite time is a sure event. 412 PHILIP L. SMITH FIG. 1. Model for simple and GoNo Go RT. The time-dependent information function +(t) is perturbed by white noise W(t) and accumulated. A response is emitted when the accumulated decision stage activation X(t) exceeds a criterion a. to accumulation (e.g., Link, 1975, 1992; Ratcliff, 1978; Smith 6 Vickers, 1988). The accumulation process then operates on the difference function, +(t)&c, rather than on the information function directly. The two panels of Fig. 2 show two ways in which the accumulation process may operate in this situation. In the model of Fig. 2a, the difference function +(t)&c is split into positive and negative halves. This is represented in the figure by a pair of opposite-signed, halfwave rectifiers, whose outputs are [+(t)&c] + =max(+(t)&c, 0) and [+(t)&c] & =&min(+(t)&c, 0), respectively. These functions are perturbed by independent sources of noise, W 1(t) and W 2(t), to yield two independent accumulation functions, X 1(t) and X 2(t). These functions race one another, with the response being determined by which of the two criteria, a 1 or a 2 , is first exceeded. Derivation of response time statistics for this model involves a pair of one-sided first passage time problems.2 In a variant of this scheme, a single noisy information function, +(t)+W(t), is split into positive and negative halves, again by the device of comparing it to a referent. The statistical properties of the accumulation function are then induced by the action of the pair of rectifiers operating on the noisy information function directly. Models of this kind belong to the class that Smith and Vickers (1988) called parallel integratorsa class which includes the recruitment model (La Berge, 1962), the Poisson parallel counter model (Pike, 1966, 1968; Townsend 6 Ashby, 1983), and the continuous state space accumulator model (Smith 6 Vickers, 1988, 1989; Vickers, 1970, 1979). In the model of Fig. 2b, the noise-perturbed difference function +(t)+W(t)&c drives a single accumulation process, X(t), whose average rate of change is either positive or negative, depending on the stimulus presented. To make a decision, the subject sets a pair of criteria, a 1 and a 2 , on the accumulated information axis and responds according to which of the events X(t)a 1 or X(t)a 2 first occurs. Derivation of response time statistics for this model requires solution of a two-sided 2 The model of Fig. 2a makes the somewhat artificial assumption that noise is localized to stages of processing occurring after rectification. Relaxation of this assumption typically leads to models in which the noise across channels is correlated. Ways to deal with correlated noise in diffusion process models are described in the last section of this article. First passage time problems for such processes are not particularly tractable, except in special cases. STOCHASTIC DYNAMIC MODELS 413 FIG. 2. Models for choice RT. (a) Parallel channels model. The function +(t) is compared to a sensory referent c to produce a signed information function +(t)&c. This is split into positive and negative parts [+(t)&c] + and [+(t)&c] & by the action of a pair of half-wave rectifiers. The rectified functions are perturbed by independent sources of white noise, W 1(t) and W 2(t), and accumulated as separate evidence totals, X 1(t) and X 2(t). The response is r 1 or r 2 , depending on which of the events X 1(t)a 1 or X 2(t)a 2 first occurs. (b) Two-barrier single channel model. The signed information function +(t)&c is perturbed by a single noise source W(t) and accumulated as a single signed total X(t). The response is r 1 or r 2 depending on whether X(t)a 1 or X(t)a 2 occurs first. first passage time problem. Models that conform to this scheme include the various random walks (Ashby, 1981; Edwards, 1965; Heath, 1981; Laming, 1968; Link, 1975, 1992; Link 6 Heath, 1975; Stone, 1960) and their continuous-time counterparts, the diffusion process models (Busemeyer 6 Townsend, 1992, 1993; Heath, 1992). The parallel diffusion process model of Ratcliff (1978, 1981) combines the characteristics of both of the models in Fig. 2. Models of simple decisions that have been proposed in the literature typically have made two important simplifying assumptions about the accumulation process X(t). These are that it is (a) time-homogeneous and (b) an independent-increments process. A time-homogeneous process is one whose transition distribution satisfies the relation P[X(t)x | X({)= y]=P[X(t&{)x | X(0)= y], (4) 414 PHILIP L. SMITH for all {<t. In such a process, the conditional probability of a transition from state y to state x depends only on the interval t&{ and not on the time at which the transition occurs. An independent-increments process is one whose transition distribution satisfies the relation P[X(t)x | X({)= y]=P[X(t)x& y | X({)=0]. (5) In such a process, the conditional probability of a transition from state y to state x depends only on the difference between the states x& y and not on the initial state y. In other words, an independent-increments process is spatially homogeneous (e.g., Chung 6 Williams, 1983, p. 12). The assumption that X(t) is time-homogeneous is equivalent to the assumptions that (a) the information function +(t) in Figs. 1 and 2 is constant and (b) the noise function W(t) is a (strictly) stationary stochastic process; that is, none of its statistics change with time. The assumption that X(t) is an independent-increments process is entailed by the frequently made assumption that the decision process is a perfect integrator; that is, accumulation of information is neither bounded nor subject to decay. While both of these assumptions may be defensible, at least in some settings, there are other occasions in which assumptions of greater dynamical complexity are warranted. Frequently, a modeler may wish to assume that the characteristics of the stimulus information change with time. Such changes may arise for a number of reasons. Most obviously, stimulus information may be available only for a limited time, as occurs under conditions of tachistoscopic presentation. Formal models for performance in this situation, which embody a single, abrupt, stimulus-driven change in the accumulation process, were described by Heath (1981) and Ratcliff (1980). More generally, continuous, smooth changes in the accumulation process may arise internally, through the action of mechanisms that encode the stimulus. Following de Lange (1952, 1954, 1958), the psychophysical properties of the early stages of visual processing are often modeled as a set of low-pass linear filters (e.g., Busey 6 Loftus, 1994; Sperling 6 Sondhi, 1968; Watson, 1986). Within this framework, timedependent variations in the statistics of the stimulus representation are assumed to arise because of the phase and frequency response characteristics of the filters. A model combining linearly filtered stimulus encoding with a stochastic accumulation process was first suggested by Heath (1992). Smooth changes in the dynamics of stimulus encoding may also arise if the underlying mechanism is selectively sensitive to changes in the stimulus input rather than to steady state intensity levels. Psychophysically transient systems of this kind have been identified in both vision (e.g., Legge, 1978; Tolhurst, 1975a, 1975b) and audition (e.g., Abeles 6 Goldstein, 1972; Gerstein, Butler, 6 Erulkar, 1968) and their properties used to explain performance in models of RT by Burbeck and Luce (1981), Burbeck (1985), and Rouder (2000). Finally, such changes may occur because of shifts in a subject's attention during the course of a trial. Depending on the nature of the task, such changes may be thought of as a set of discrete, punctate transitions, as proposed in the decision model of Diederich (1997) or as smoothly varying changes in the rate of accumulation, as suggested by the attention gating STOCHASTIC DYNAMIC MODELS 415 model of Sperling and colleagues (Reeves 6 Sperling, 1986; Sperling 6 Weichselgartner, 1995). Regardless of how such variations are produced, however, they imply timeinhomogeneity in the accumulation process X(t). Violations of the independent-increments assumption occur when there is state dependence in the accumulation processindeed, this is the substance of the definition in (5). State dependencies of this kind arise when the signal-to-noise ratio (SNR) of the accumulation process is bounded, as occurs in imperfect or leaky accumulators. In leaky accumulators, the growth of activation in the decision system caused by the presence of a stimulus is opposed by a tendency for activation to decay at a rate that is proportional to its current level. Models with stochastic, SNR-bounded accumulation have been considered by Busemeyer and Townsend (1992, 1993), Diederich (1995), Pacut (1977, 1980), and Rudd and Brown (1997). A model combining all of the preceding dynamic attributes, namely, linearly filtered stimulus encoding, change and level sensitive mechanisms, and leaky stochastic accumulation, was proposed by Smith (1995, 1998a). The preceding paragraphs have considered a number of examples from the areas of perception and performance in which violations of either spatial homogeneity or temporal homogeneity, or both, arise. Many of these considerations also apply, either implicitly or explicitly, to mathematical models of memory. Although the simplest versions of most memory models are static (e.g., Gillund 6 Shiffrin, 1984; Humphreys, Bain, 6 Pike, 1989; Murdock, 1982) and make predictions that can be characterized using signal detection theory alone, the theoretical framework that accompanies them is usually rich enough to support more complex dynamic predictions if they are required (e.g., Ratcliff, 1978). Dynamic behavior of this kind arises naturally, for example, with the addition of assumptions that specify the time course of the interaction between a probe stimulus and the memory system. Indeed, such dynamics are an explicit part of some network models of memory, such as interactive activation models (Coltheart 6 Rastle, 1994; McClelland 6 Rumelhart, 1981). In general, such models postulate time-dependent and asymptotically bounded information accrual, which results in accumulation functions that are both spatially and temporally inhomogeneous. STOCHASTIC INFORMATION ACCRUAL AS A DIFFUSION PROCESS This article is concerned with diffusion process models of information accrual; that is, continuous-time, continuous sample path Markov processes. Although there are other ways in which information accrual in continuous time may be modeled the most important among them from a biological perspective being point-process models (e.g., Diederich, 1995; Luce 6 Green, 1972; Hildreth, 1979; McGill, 1963; Rudd, 1996; Schwartz, 1992; Smith 6 Van Zandt, in press; Townsend 6 Ashby, 1983)diffusion processes are of particular interest because they include an important class of continuous-time Gaussian processes, which provide a natural, dynamic generalization of the Gaussian signal detection models that form the mainstay of much psychological theorizing. Furthermore, Gaussian approximations may be appropriate to model molar accumulation behavior even when it is assumed that the fine-grained behavior of the system can be described as a point process, as, for 416 PHILIP L. SMITH example, when modeling the number of active fibers in a neural relay. In this situation, a possible model for the accumulated impact of successive spike discharges on a post-synaptic membrane potential is as Poisson shot noise, which becomes approximately Gaussian at high intensities of the incident point process (e.g., Papoulis, 1991, pp. 629635). 3 Variants of this idea appear in Link (1992), McGill (1967), Rudd (1996), and Smith (1998b), among others. The rest of the article considers in turn the two problems described previously: first, how to characterize the stochastic properties of the accumulation process X(t); second, how to solve the one-sided and two-sided first passage time problems for X(t) when it is spatially and temporally inhomogeneous. The classical approach to first passage time problems (e.g., Cox 6 Miller, 1965, Ch. 5; Feller, 1971, Ch. 10) treats them as boundary value problems in the theory of partial differential equations and yields tractable solutions only in the time-homogeneous case. A second approach that has been applied to cognitive problems is to approximate the process using a finite-state Markov chain and to compute the first passage time statistics for the approximating process using spectral methods, as was done by Busemeyer and Townsend (1992, 1993) and Diederich (1995, 1997). This method yields computationally efficient solutions for spatially homogeneous and inhomogeneous problems, but only in the time homogeneous or piecewise time homogeneous cases. In contrast, recent developments in the applied probability literature have led to integral equation representations of the solution of the one-sided and two-sided first passage time problems that can often yield tractable results even in the presence of time inhomogeneity. These methods, which provide high-accuracy numerical approximations to the first passage time distributions, are of particular relevance to the cognitive modeler and form the basis of the procedures described here. Classically, F(x, t | y, {), the transition distribution of a diffusion process X(t), is shown to satisfy a pair of partial differential equations known as the backward and forward (or FokkerPlanck) equations. Solution of the relevant equation in the presence of an initial condition yields the transition distribution; solution in the presence of appropriate boundary conditions yields the first passage time distribution. These methods are described in Cox and Miller (1965, Ch. 5), Feller (1971, Ch. 10), and Karlin and Taylor (1981, Ch. 15). Applications of these methods in a cognitive setting are described by Ratcliff (1978, 1980) and Smith (1990). A more modern approach to the study of diffusions is via stochastic differential equations (SDEs). Such an approach has the advantage of providing a characterization of the process that is more direct and intuitive than that afforded by partial differential equations. Its disadvantage is that a rigorous treatment of SDEs requires the analytic framework of measure theory and is thus unsuitable for a brief, self-contained presentation such as this. Fortunately, however, the SDE approach to the important class of 3 Pooling of activation across positive and negative channels of the kind found in the human visual system may be modeled as the difference of a pair of Poisson shot noise processes. At high intensities, this process will also approximate a Gaussian process. Such processes may provide a theoretical link between the underlying neural mechanisms and the diffusion process models described in this article, but no attempt is made here to develop these properties in detail. The interested reader is referred to Rudd (1996) and Rudd and Brown (1997) for highly developed models which link patterns of neural firing, diffusion processes, and early vision. STOCHASTIC DYNAMIC MODELS 417 Gaussian processes can be motivated heuristically in a fairly simple way, by considering the limits of sums. This is the approach I adopt here. Useful introductions to the theory of SDEs may be found in Arnold (1974), Bhattacharya and Waymire (1990), Gardiner (1985), and Karlin and Taylor (1981). Rigorous measure-theoretic treatments are contained in Ethier and Kurtz (1986), Karatzas and Shreve (1991), Protter (1990), and Revus and Yor (1994). We wish to characterize a class of stochastic accumulation processes that is described by the SDE dX(t)=+(X(t), t) dt+_(X(t), t) dB(t). (6) In this equation, dX(t) is the random change in the process X(t) that occurs in a small time interval dt. We seek to interpret this equation as the limit, in some suitable sense, of the discrete time difference equation X(t+2t)&X(t)=+(X(t), t) 2t+_(X(t), t) 2B(t), (7) in which 2B(t)=B(t+2t)&B(t) is a small Gaussian-distributed perturbation of order - 2t. Symbolically, we write 2B(t)t- 2t. The requirement that 2B(t) approaches zero more slowly than 2t is designed to ensure that the stochastic variation in X(t) is preserved in the limit. We ascribe the following properties to 2B(t): E[2B(t)]=0 E[2B 2(t)]=2t (8) Cov[2B(t+2t) 2B(t)]=0. Equation (8) stipulates that (a) 2B(t) is a zero-mean Gaussian process; (b) its variance in any interval 2t equals the length of the interval, and (c) the increments in successive nonoverlapping intervals are independent. The physical intuition that underlies the description of X(t) in (7) is as follows: The change 2X(t) in any interval 2t consists of two parts, one deterministic, the other stochastic. The magnitudes of the two parts depend on the functions +(X(t), t) and _ 2(X(t), t), respectively. Both parts may depend on state x and on time, t, in which case, X(t) is spatially and temporally inhomogeneous. Because 2B(t) is Gaussian, the increment 2X(t) will be Gaussian with mean +(X(t), t) 2t and variance _ 2(X(t), t) 2t. The resulting process X(t) may be thought of as a realization of the trajectory in state space of a system whose underlying smooth dynamics are prescribed by the function +(X(t), t), but which is subject to repeated shocks or statistical perturbations. For a given set of initial conditions, the actual trajectory of the system will depend jointly on the factors that determine its smooth dynamics and on the unique realization of the sequence of perturbations 2B(t i ), i=1, 2, ... . 418 PHILIP L. SMITH FIG. 3. Solution to Eq. (6) as a diffusion process X(t), with initial condition X(t 0 )=x 0 . The value of X(t) is a random variable with distribution function F(x, t | x 0 , t 0 ), transition density f (x, t | x 0 , t 0 ), and mean m(t; t 0 ). The transposed normal curve suggests the properties of the transition density function graphically. Within this framework, (6) may be thought of as describing the trajectory of X(t) when the interval between successive shocks becomes very small, as suggested by the sample paths shown in Fig. 3. To interpret (6) as the limit of (7) as 2t goes to zero requires us to give meaning to quantities of the form W(t)=lim 2 Ä 0 2B(t)2t. However, the stipulation that 2B(t)t- 2t precludes convergence in any ordinary sense (e.g., Karlin 6 Taylor, 1981, p. 341). Indeed, the prescriptions contained in (8b) and (8c) imply that Cov[W(t) W({)]=$(t&{), the Dirac delta function which does not exist as a function in any usual sense. Nevertheless, a generalized function interpretation may be given to W(t). Recalling (a) that the power spectrum of a stochastic process is the Fourier transform of its autocorrelation function and (b) that the Fourier transform of the Dirac delta function is a constant, W(t) may be interpreted as a Gaussian process whose power spectrum is constant; that is, as one in which there is equal power at all frequencies. Such a process is known as white noise (Wong 6 Hajek, 1985, pp. 109115). As a process, white noise cannot exist in any real physical sense. Nevertheless, it provides a convenient and widely used approximation when describing the dynamics of physical systems that are additively perturbed by broad-spectrum Gaussian noise. Indeed, as noted by Karatzas and Shreve (1991), the development of a rigorous stochastic calculus by Ito^ in the 1940s was motivated by the need to understand systems of this kind. Here we use only the minimum of this apparatus needed to provide a characterization of the distributional properties of X(t). The interested reader is referred to the previous references for further details. Under appropriate smoothness conditions on the functions +(x, t) and _(x, t), the solution X(t) to (6) will be a continuous sample path Markov process. Given STOCHASTIC DYNAMIC MODELS 419 certain technical restrictions on the process X(t), a sufficient condition for it to have continuous sample paths is that it satisfy the Dynkin condition 4 lim hÄ0 1 P[( |X(t+h)&X(t)| >=) | X(t)=x]=0 h for all =>0. Loosely, the probability of large jumps in X(t) in any interval goes to zero with the size of the interval. In particular, this condition excludes jump processes, the canonical example of which is the Poisson process. 5 The resulting process is fully characterized by its infinitesimal moments, +(x, t) and _ 2(x, t), which are defined as: 1 E[X(t+h)&X(t) | X(t)=x] hÄ0 h +(x, t)= lim (9) 1 _ (x, t)= lim E[[X(t+h)&X(t)] 2 | X(t)=x]. hÄ0 h 2 These functions, which are known as the drift and diffusion coefficients of the process, respectively, are the infinitesimal (rate) equivalents of the mean and variance of the increment process in (7). As written, the SDE in (6) includes both linear and nonlinear processes of arbitrary order. Here we confine ourselves to the general, linear, first-order SDE dX(t)=[+(t)+b(t) X(t)] dt+[_(t)+c(t) X(t)] dB(t), (10) in which +(t), b(t), _(t), and c(t) are given continuous functions of timewhich includes the possibility that they may be constant or zero. Of the various special cases of (10), the most important is that in which c(t) is zero, i.e., dX(t)=[+(t)+b(t) X(t)] dt+_(t) dB(t). (11) In this equation the diffusion term is either constant or may depend on time but is independent of state. The importance of this equation is that its solutions are Gaussian processes and, unlike the solutions of the more general equation (10), such solutions may be obtained by elementary methods. In general, SDEs are 4 The requirements are that X(t) is right continuous and possesses left limits. A process is right continuous if lim t a { X(t)=X({) for all {. It possesses left limits if lim t A { X(t) exists for all {>0. A strong Markov process with these properties that satisfies the Dynkin condition is a diffusion process. Most processes of interest in applications either satisfy these conditions or can be realized as processes of this kind. In particular, any strong Markov process which is continuous in probability (i.e., for which lim t Ä { P[ |X(t)&X({)| >=]=0 for any {0, =>0) has an equivalent version which satisfies these conditions. A criterion for the Dynkin condition to hold is that the infinitesimal moment condition lim h Ä 0 E[ |X(t+h)&X(t)| p | X(t)=x]h=0 holds uniformly in x for some p>2 (Karlin 6 Taylor, 1981, p. 165). 5 Processes consisting of the superposition of a continuous process and a jump process are called Levy processes. Such processes are discussed by Protter (1990). 420 PHILIP L. SMITH usually written in the differential form used in (10) and (11) rather than the more familiar form involving derivatives because of the difficulty in assigning meaning to expressions of the form dXdt when the function X(t) is not of finite variation. 6 Before proceeding, we distinguish two further special cases of (11). These are dX(t)=+(t) dt+_ dB(t) (12) dX(t)=[+(t)&#X(t)] dt+_ dB(t), (13) and which are obtained by setting _(t)=_ constant and b(t)=0 or b(t)=&#. Equations (12) and (13) generate, respectively, the two diffusion processes that historically have been of the greatest importance in both theory and applications, namely the Brownian motion (or Wiener) process and the OrnsteinUhlenbeck (OU) process. 7 These processes are also the ones that have found the greatest application in sensory and cognitive modeling (e.g., Rudd 6 Brown, 1997; Busemeyer 6 Townsend, 1992, 1993; Diederich 6 Busemeyer, 1995; Emerson, 1970; Heath, 1992; Ratcliff, 1978, 1981; Reed, 1973; Smith, 1995, 1998a). They may be interpreted as describing the dynamics of a perfect and a leaky integrator, respectively. Detailed accounts of the use of the OU process to model leaky integration may be found in Busemeyer and Townsend (1992) and Smith (1995). We attempt to give meaning to (11) using integration by parts. Proceeding by analogy with the deterministic case (Karlin 6 Taylor, 1981, pp. 345346), we introduce the integrating factor exp[& t b({) d{] and write (11) in the form e & t b({) d{ [dX(t)&b(t) X(t) dt]=e & t b({) d{ [+(t) dt+_(t) dB(t)]. The left-hand side of this equation can be recognized as an exact differential: d[X(t) e & t b({) d{ ]=e & t b({) d{ [+(t) dt+_(t) dB(t)]. Both sides of the equation may therefore be integrated from 0 to t and the result rearranged to yield X(t)= | t 0 t e { b(s) ds+({) d{+ | t t e { b(s) ds_({) dB({), (14) 0 where, as stipulated previously, the initial condition P[X(0)=0]=1 has been assumed. Equation (14) may be interpreted as the output of a linear system with 6 The variation of a process is defined in the following way: Let > (m) =[t 0 , t 1 , ..., t m ] be a partition of the compact interval [0, t]. The total variation of X(t) on [0, t] is defined as sup m i=1 |X(t i )&X(t i&1 )|, where the supremum is taken over all possible subdivisions of [0, t] with m arbitrarily large. A process is said to be of finite variation if its total variation is finite. A fundamental property of the sample paths of diffusion processes is that, with probability one, their total variation on compact intervals is infinite. 7 For historical reasons, SDEs of the form (13) are known as Langevin equations (see, e.g., Gardiner, 1985; van Kampen, 1992). 421 STOCHASTIC DYNAMIC MODELS impulse response function b(t) (cf. Norman, 1981) and input +(t)+_(t) dB(t). That is, the output is the superposition of outputs obtained by independently supplying as inputs the deterministic function +(t) and the temporally modulated white noise process _(t) dB(t). To interpret the stochastic integral on the right of (14) we use a second integration by parts to get | t t e { b(s) ds_({) dB({)=_(t) B(t)+ 0 | t t e { b(s) ds[_({) b({)&_$({)] B({) d{. (15) 0 In this form, the integral on the right of (15) may be seen to be a linear function of the Brownian motion process B(t), which is the canonical independent-increments diffusion process (Karatzas 6 Shreve, 1991). By virtue of the independent-increments property, the functional central limit theorem ensures that B(t) is Gaussian (Bhattacharya 6 Waymire, 1990, pp. 2024) and thus can be fully characterized by its first two moments, namely, E[B(t)]=0 Cov[B(t) B({)]=min(t, {), and, in particular, Var[B(t)]=t (cf. (10)). As a linear functional of a Gaussian process, the stochastic integral in (15) will also be Gaussian, and X(t), as a sum of deterministic and stochastic parts, will inherit this property also. We seek to ascertain the mean, variance, and transition distribution of X(t). Before proceeding, we observe that the special case b(t)=&#, _(t)=_ in (14) and (15) yields a linear functional representation of the OU process (Bhattacharya 6 Waymire, 1990, p. 581; Karlin 6 Taylor, 1981, p. 345) X(t)= | t _ e &#(t&{)+({) d{+_ B(t)&# 0 | t & e &#(t&{)B({) d{ , 0 (16) and, trivially, when #=0, of the (time-inhomogeneous) Brownian motion process X(t)= | t +({) d{+_B(t). (17) 0 To motivate subsequent formal manipulations, we seek to ascertain the distributional properties of the process X(t) in (16) heuristically, as the limit of sums. An analogous, but more laborious, argument could be constructed for the more general process (15) but would have little additional heuristic value. To this end, we consider a special case of the difference equation (7) 2X i =+~ &#~X i&1 +_~W i , 422 PHILIP L. SMITH in which 0<#~ <1 and W i , i=1, 2, ..., is a sequence of independent and identically distributed Gaussian random variables with E(W i )=0 and Var(W i )=1. This equation may be written X i =+~ +(1&#~ ) X i&1 +_~W i and solved recursively to yield n&1 n&1 X n = : (1&#~ ) i +~ + : (1&#~ ) i _~W i . i=0 (18) i=0 To obtain moments when the interval between successive increments is arbitrarily small, we arrange the passage to the limit such that n=t2, _ 2 2=_~ 2, # 2=#~, and +2=+~. With this substitution, (18) becomes n&1 n&1 X n =2+ : (1&2#) i +- 2_ : (1&2#) i W i . i=0 (19) i=0 Taking expectations in (19), and using the expression for the partial sum of a geometric series (convergence of which is guaranteed by the condition on #~ ), yields the expected value n&1 E[X n ]=2+ : (1&2#) i i=0 1&(1&2#) n 1&(1&2#) _ & 1&(1&#tn) =2+ _ 2# & . =2+ n The standard identity lim n Ä (1+xn) n =e x can then be applied to give + E[X t ]= lim E[X n ]= [1&e &#t ]. 2Ä0 # (20) The variance is obtained in a similar way, Var(X n )=E _{ n&1 - 2 _ : (1&2#) i W i i=0 n&1 =2_ 2 : (1&2#) 2i E[W 2i ], i=0 2 =& 423 STOCHASTIC DYNAMIC MODELS where the vanishing of the cross-product terms occurs by virtue of the fact that Cov(W i W j )=0, i{ j. As before, the series may be summed to give 1&(1&2#) 2n 1&(1&2#) 2 _ & 1&[(1&#tn) ] =_ 2 _ 22#&2 # & . Var(X n )=_ 2 2 n 2 2 2 2 By a similar passage to the limit to that which yielded (20) Var(X t )= lim Var(X n )=_ 2 2Ä0 { 1&e &2#t . 2# = (21) For the special case in which #=0, a similar, but simpler, calculation gives E[X t ]=+t, Var(X t )=_ 2t. (22) Equations (20)(22) give the mean and variance of the OU process and the Brownian motion process as limits of sums for the time-homogeneous case, +(t)=+. From these pairs of expressions, an important difference between the two processes is apparent, whose psychological significance was first discussed by Busemeyer and Townsend (1992). For the Brownian motion process, the time-dependent SNR is E 2[X t ]Var(X t )=(+ 2_ 2 ) t; that is, the separation between signal and noise grows unboundedly with t. In contrast, for the OU process E 2[X t ]Var(X t ) Ä 2+ 2(#_ 2 ); that is, the SNR grows to an asymptotic limit. This difference between the two processes is a reflection of their characterizations as perfect and leaky integrators, respectively. Implicit in the formal manipulations used to obtain (20) and (21) as a limit of sums was the assumption that an exchange of limit and integral of the form lim n Ä E[X n ] =E[lim n Ä X n ] was possible. 8 A rigorous justification for these manipulations is provided by the Ito^ calculus, which also yields a way to obtain the moments of the limit processes in (14) and (16) directly. To show the parallels between the results obtained using these methods and those obtained from the limit of sums, we consider the solution to the SDE (13) for the time-homogeneous OU process (+(t)=+): X(t)= | t 0 e &#(t&{)+ d{+_ | t e &#(t&{) dB({). (23) 0 An important property of the stochastic integral on the right in (23) is that, when interpreted in an Ito^ sense, it preserves the martingale property of the integrator 8 The preceding construction suggests, but does not prove, that the sample paths of the limit process X t are continuous. To obtain a continuous-sample path process with the requisite properties requires a more elaborate construction than the one given here. Details may be found in Karlin and Taylor (1975, pp. 371378) or Karatzas and Shreve (1991, pp. 5659). 424 PHILIP L. SMITH B(t). 9 In general, a stochastic process Z(t) is said to be a martingale if E[ |Z(t)| ] < for all t, and E[Z(t) | F{ ]=Z({), {<t, where Ft is the sigma-field generated by the process [Z({), 0{<t]. Loosely, Ft represents all of the information that can be observed about Z(t) up to the time t. Specifically, for the Brownian motion process, B(t), the martingale property implies that E[B(t) | B(0)=0]=0. Because the stochastic integral preserves the martingale character of B(t), _ E _ | t } & e &#(t&{) dB({) B(0)=0 =0, 0 and therefore E[X(t)]= | t 0 + e &#(t&{)+ d{= (1&e &#t ). # In other words, the expected value of X(t) is the output of a linear system with impulse response function &# to a constant input +. This result agrees with that obtained from a limit of sums in (19). To obtain the variance of X(t) we use the fundamental Ito^ isometry (Chung 6 Williams, 1983, p. 27; Karatzas 6 Shreve, 1991, pp. 137138) E _{| 2 t X({) dB({) 0 =& = | t E[X 2({)] d{. (24) 0 The mean-square (L 2-norm) convergence of this equation is a crucial step in establishing the existence of the stochastic integral in a rigorous way. The equation itself may be derived by considering the limits of sums and from the properties dB(t)t- dt, E[dB 2(t)]=dt, and Cov(dB(t) dB({))=0. Applying (24) to the stochastic integral in (23) yields Var(X(t))=_ 2E =_ 2 = | _{| t t 2 e &#(t&{) dB({) 0 =& e &2#(t&{) d{ 0 _2 [1&e &2#t ], 2# (25) 9 Stochastic integrals may be defined with respect to integrators other than Brownian motion. The class of integrators for which the Ito integral is defined are called semimartingales. Such a process can be decomposed into a sum of a finite variation process and a continuous local martingale. The latter are processes which may fail globally to be martingales (e.g., because they are unbounded), but which exhibit the martingale property at each of an increasing sequence of stopping times. The Ito^ integral with respect to such processes preserves the local martingale property of the integrator. See Karatzas and Shreve (1991) or Protter (1990), for details. 425 STOCHASTIC DYNAMIC MODELS again in agreement with the result obtained in (21) as a limit of sums. The preceding method may be applied straightforwardly to obtain the moments for the general Gaussian process in (14). From a knowledge of the first two moments of X(t) and the fact that it is Gaussian, the transition distributionor, equivalently, the transition densityof X(t) may be written down immediately. The transition density f (x, t | y, {) is defined as f (x, t | y, {)=(ddx) F(x, t | y, {). (26) The properties of this density are suggested graphically in Fig. 3. For later reference we write down the transition densities for the time-inhomogeneous Brownian motion and OU processes. To this end, we let the function m X (t; {) denote the time-dependent mean for a process starting at time { evaluated at time t. The transition density for the inhomogeneous Brownian motion process in (12) is f (x, t | y, {)= 1 - 2?_ 2(t&{) _ exp & (x& y&m X (t; {)) 2 2_ 2(t&{) & (27) with m X (t; {)= | t +(s) ds. { The corresponding density for the inhomogeneous OU process in (13) is f (x, t | y, {)= # ?_ [1&exp[&2#(t&{)]] &#[x& y exp[&#(t&{)]&m (t; {)] _exp { _ [1&exp[&2#(t&{)]] = , 2 2 X 2 (28) with m X (t; {)= | t e &#(t&s)+(s) ds. { These densities will be used subsequently in the derivation of solutions to the onesided and two-sided first passage time problems for X(t). The general, linear SDE in (10), despite its superficial similarity to the Gaussian equation (11), cannot be dealt with in the same elementary way. Attempts to do so result in integrals containing terms of the form X(t) dB(t), in which both the integrator and the integrand are stochastic, and these cannot be reduced to a Gaussian functional by the device of twice integrating by parts. Rather, the key to solving equations of the form (10) is the Ito^ transformation formula, which is the stochastic counterpart of the 426 PHILIP L. SMITH change of variables formula from ordinary calculus. Specifically, let Z(t) be a stochastic process that satisfies the SDE dZ=+(Z(t), t) dt+_(Z(t), t) dB(t), and let f (z, t) be a continuous function of two variables which is twice-differentiable in its first argument and once-differentiable in its second. Let f $z and f "zz denote, respectively, the first and second partial derivatives of f with respect to z and let f $t denote the first partial derivative with respect to time. Then the transformed process f(Z(t), t) satisfies the SDE df (Z(t), t)=[ f $z (Z(t), t) +(Z(t), t)+ f $t(Z(t), t)+ 12 f "zz (Z(t), t) _ 2(Z(t), t)] dt + f $z (Z(t), t) _(Z(t), t) dB(t). (29) Formally, (29) is derived by expanding f (Z(t), t) in a truncated Taylor series around (Z(t), t) and discarding terms of order dB(t) dt, (dt) 2, and higher. The key to performing this expansion in a stochastic setting is the fact that dB 2(t)tdt. This yields f (Z(t+dt), t+dt)= f (Z(t), t)+ f $z (Z(t), t) dZ(t)+ f $t(Z(t), t) dt + 12 f "zz (Z(t), t) d 2Z(t) or df (Z(t), t)= f $z (Z(t), t)[+(Z(t), t) dt+_(Z(t), t) dB(t)] + f $t(Z(t), t) dt+ 12 f "zz (Z(t), t) _ 2(Z(t), t) dt, because d 2Z(t)t_ 2(Z(t), t) dt. This result differs from the one that would be obtained if Z(t) were deterministic by the addition of the term f "zz (Z(t), t) _ 2(Z(t), t) dt2, reflecting the inclusion in the expansion of the term d 2B(t). The presence of this term means that the transformation rules for SDEs differ from those of ordinary calculus. In the Appendix the Ito^ transformation rule is used to obtain the solution of the general linear SDE (10) and of the related homogeneous equation dX(t)=b(t) X(t) dt+c(t) X(t) dB(t). (30) Specifically, let X(0) be the (random) initial value of X(t) and define the auxiliary function U(t)=exp _| t b({) d{+ 0 | t 0 c({) dB({)& 12 | t 0 & c 2({) d{ . (31) The solution of the homogeneous equation (30) may then be written X(t)=X(0) U(t). (32) 427 STOCHASTIC DYNAMIC MODELS The solution of the general equation (10) is { X(t)=U(t) X(0)+ | t 0 1 [+({) d{+_({) dB({)&c({) _({) d{] . U({) = (33) From a modeling perspective, it is important to know why one would choose the representation (10) instead of (11) (or vice versa). The SDE (11), in which the diffusion term may depend on time, but not on state, is appropriate when the noise in the decision stage is exogenous or input noisethat is, noise arising because of statistical variation in the stimulus encoding processas distinct from noise intrinsic to the decision stage itself. The possibility for time-inhomogeneity in the drift and diffusion terms expressed by (11) allows for various forms of nonstationarity in the encoding process, whose mean and variance may vary jointly as a function of time. One possibility of theoretical interest that can be represented in this way is Poissonlike noise (cf. Link, 1992), which is obtained by setting _(t)=- +(t) in (11). With this constraint, the stimulus-driven increments to X(t) form a Gaussian process whose time-dependent mean and variance are equal. Another possibility is Weberlike noise, which is obtained by setting _(t)=+(t). This represents a Gaussian process in which the standard deviation (rather than the variance) of the increment process grows in proportion to its mean. An alternative perspective to that which relegates all of the noise in the system to the input side is to view the decision stage as analogous to an audio amplifier, in which noise grows in proportion to signal strength or, in this case, in proportion to the level of activation in the decision stage. In these circumstances, the dynamics of the system are likely to be represented more appropriately by a variant of (10), in which the variance of the increments to X(t) is determined endogenously rather than exogenously, by the instantaneous value of X(t). The homogeneous equation (30) describes a system of this kind, in which the noise is purely endogenous. The general linear equation (10) that is solved by (33) represents a system in which exogenous, stimulus-driven noise and endogenous decision stage-noise both contribute to overall system variability. Although the representation in (33) is attractive in its generality, a number of complications attend its use in cognitive models. First, as noted previously, the process X(t) in (33) is in general not Gaussian. Higher moments of X(t) may be obtained using a method described in Gardiner (1985, pp. 112113), but no general closed form expression for the transition density of the process appears to exist. This limitation is unimportant if the process X(t) is used to represent a timedependent signal detection model, in which the sample size is determined by factors external to the accumulation process, as knowledge of the moments suffices to determine the required statistics. However, it is relevant to models in which a first passage time problem must be solved, which require an explicit expression for the transition density, as described in the following section. One important special case for which an explicit representation of the transition density exists is the homogeneous process (32). Although the process X(t) in this equation is not Gaussian, it is immediate from (31) and (32) and the previous characterization of independentincrements processes that a Gaussian process is obtained by the transformation 428 PHILIP L. SMITH log[X(t)X(0)]=log U(t). Indeed, from the martingale property of the stochastic integral and the Ito^ isometry (24), the mean and variance of log[X(t)X(0)] are m X (t; 0)= | t 0 b({) d{& 12 | t c 2({) d{ 0 and v X (t; 0)= | t c 2({) d{, 0 respectively, where the notation is that of (27) and (28). This means that the process X(t) is lognormal (i.e., a random variable whose logarithm is Gaussian) with transition density f(x, t | y, {)= 1 { exp & x - 2?v X (t; {) [log(xy)&m X (t; {)] 2 . 2v X (t; {) = (34) A second complication involves the existence of competing interpretations of the stochastic integral in the solution to the SDE (10). The difficulty in finding an explicit representation for the process X(t) that solves this equation arises from the problem of giving meaning to stochastic integrals such as Z(t)= | t X({) dB({), 0 in which both the integrator and the integrand are stochastic. The martingale property that was used to derive moments of the processes X(t) in (23) and (32) arises from the presumption that all stochastic integrals are interpreted in a particular way. According to this Ito^ interpretation, the integral is constructed as a limit of sums in the following manner: Let 6 (m) =[t 0 , t 1 , ..., t m ] be a partition of the interval [0, t] such that 0=t 0 <t 1 < } } } <t m =t. The Ito^ integral is defined to be the limit in mean square as the mesh of 6 (m) goes to zero of the sum | m&1 t X(t) dB(t)= 0 lim-m.s. max(ti+1 &ti ) Ä 0 : X(t i )[B(t i+1 )&B(t i )]. i=0 A crucial part of this definition is the fact that the integrand in the interval (t i , t i+1 ] is defined to equal X(t i ), its value at the left end point of the interval. In other words, the integrand is nonanticipating. The martingale-preserving property of the stochastic integral requires that it be defined in this way. An alternative definition of the stochastic integral was proposed by Stratonovich (e.g., Gardiner, 1985; Karatzas 6 Shreve, 1990; Karlin 6 Taylor, 1981) in which it is defined instead to be the limit of sums of the form | t 0 m&1 X(t) b dB(t)= lim-m.s. max(ti+1 &ti ) Ä 0 : i=0 1 2 [X(t i+1 )+X(t i )][B(t i+1 )&B(t i )]. STOCHASTIC DYNAMIC MODELS 429 (The composition symbol `` b '' is a standard notation for the Stratonovich integral.) In this version of the integral, the value of the integrand in the interval (t i , t i+1 ] is taken to be the average of its values at the right and left endpoints. It was shown by Wong and Zakai (Karlin 6 Taylor, 1981; Wong 6 Hajek, 1985) that this version of the integral arises naturally when the integral is defined ``pathwise'', by approximating the irregular sample path of B(t) by a finite-variation Gaussian process. This is equivalent to constructing the integral as a limit of sums using as integrator a nonwhite Gaussian noise; that is, a broad spectrum noise process, B (m)(t), in which successive values of the process are not required to be delta correlated, but for which the autocorrelation function Cov[B (m)(t i+1 ) B (m)(t i )] goes to zero with the mesh of 6 (m). As discussed by Gardiner (1985) and van Kampen (1992), the Stratonovich interpretation of the stochastic integral, although it fails to preserve the martingale property of the noise process, may be more relevant to modeling real systems, in which total noise power is finite. Further, the Stratonovich integral, although it is defined for a more restricted range of integrands than is the Ito^ integral, also has the advantage that it transforms according to the familiar rules of calculus (cf. 30). Despite these differences in the definitions of the two versions of the integral, there is a systematic relationship between the Ito^ and the Stratonovich solutions of the SDE (10). Specifically, the solution of the general SDE dX(t)=+(X(t), t) dt+_(X(t), t) b dB(t), in which the stochastic integral is interpreted in the Stratonovich sense, is equivalent to the solution of the modified SDE dX(t)=[+(X(t), t)+ 12 _(X(t), t) _$x (X(t), t)] dt+_(X(t), t) dB(t), in which the integral is interpreted in the Ito^ sense. That is, the Ito^ solution and the Stratonovich solution differ by the presence or absence in the drift coefficient of a correction term that is equal to the product of half the diffusion coefficient and its derivative with respect to its space coordinate. Because of the martingale-preserving properties of the Ito^ integral, calculations are usually carried out using Ito^ integrals. The Stratonovich solution, if desired, may be obtained by adding a correction term to the original equation in this way. Details may be found in Karlin and Taylor (1981). Evidently, the behavior ascribed to the solution process X(t) may differ markedly, depending on whether the Ito^ or the Stratonovich interpretation of the stochastic integral is chosen. Unfortunately, the level of specificity required to choose between these competing interpretations in a principled way is unlikely to be present in most cognitive models, in which the noise process is hypothesized and unobservable. In contrast, no such arbitrariness exists in relation to the restricted equation (11) in which the diffusion coefficient is spatially homogeneous, because in this situation, the correction term _$x (X(t), t) is zero, and the Ito^ and Stratonovich solutions coincide. 430 PHILIP L. SMITH THE FIRST PASSAGE TIME PROBLEM FOR TIME-INHOMOGENEOUS DIFFUSION PROCESSES The Single-Barrier Case This section takes up the second of the two questions considered in this article, namely, the solution of the first passage time problem for a diffusion process in the presence of possible time inhomogeneities in the drift and diffusion coefficients of the defining SDE. The one-sided and two-sided problems will be considered in turn. We formulate the solution to this problem in a fairly general setting, in which the absorbing barrier(s) are not constrained to be fixed, but may vary smoothly with time. There are two reasons for formulating the problem in this general way. The first is that there may be a theoretical justification in some applications for assuming that the decision criteria used by subjects vary systematically with time. Dynamic variations in criteria of this kind have been proposed in various settings by a number of researchers, including Busemeyer and Rapoport (1988), Hockley and Murdock (1987), and Viviani (1979a, 1979b). The second reason for considering diffusion processes bounded by time-varying absorbing barriers is that a large class of time inhomogeneous problems can be reduced to the homogeneous case by an appropriate transformation of the state space andor the time scale. This approach, which was described by Ricciardi (1976) and Ricciardi and Sato (1983), was used in RT models described by Heath (1992) and Smith (1995). The essence of the transformation approach is that it seeks to reformulate problems involving time-inhomogeneous processes such as the OU process in (16) into equivalent problems involving time-homogeneous processes by the removal of the integrated drift term. Although we consider this approach briefly here, we prefer to work with the inhomogeneous process directly. The basic probabilistic setup is depicted in Fig. 4. As before, X(t) is a diffusion process defined on a state space R=(&, ), the real numbers, with t # [0, ), the positive real line. As discussed in the previous section, we assume that X(t) has a known transition density f (x, t | y, {). We define a(t) to be a smooth (specifically, twice-differentiable) absorbing barrier, as shown in Fig. 4, and assume that the process starts in state x 0 at time t 0 with probability one. For convenience we assume x 0 <a(t 0 ). This is the case that arises most naturally in cognitive settings, so it is the only one we consider explicitly. The solution to the converse problem x 0 >a(t 0 ) can be obtained by analogy with the one described here. We define g[a(t), t | x 0 , t 0 ]= g T (t), the first passage time density for X(t) through the timevarying absorbing boundary a(t). In this notation the subscript denoting the random variable is omitted, but the conditional dependency on the initial time, state, and absorbing barrier is made explicit. The basic tool for the analysis of such processes is a simple and intuitive renewal equation, attributed by Durbin (1971) to Fortet (1943). To obtain this equation we consider a decomposition of the probability density function of the sample path of an unconstrained process that originates at (x 0 , t 0 ) and passes through a point a(t) on the absorbing barrier at time t, as shown in Fig. 4. By definition, the probability density of sample paths passing through this point is f [a(t), t | x 0 , t 0 ]. Because STOCHASTIC DYNAMIC MODELS 431 FIG. 4. Renewal representation of the one-sided first passage time problem for a diffusion process X(t) through a time-dependent absorbing barrier a(t). Equation (35) gives a decomposition of f [a(t), t | x 0 , t 0 ], the probability density associated with a point on the absorbing barrier a(t) at time t. The time of the first barrier crossing is {, {t. After crossing at a({) the process makes a transition to a(t) in time t&{. In this figure, the symbol t does double duty to denote the time index for the process and the arbitrary point at which Eq. (35) is evaluated. x 0 <a(t 0 ), all sample paths passing through this point must have crossed the barrier at least once, at or before time t. Let the time of the first barrier crossing be {, {t. After crossing the barrier at {, the process must subsequently have made a transition from a({) to a(t) in the interval t&{. 10 The number of boundary crossings subsequent to the first is immaterial. By virtue of the (strong) Markov property of X(t) the probability densities associated with the portions of the sample path before and after { are g[a({), { | x 0 , t 0 ] and f [a(t), t | a({), {], respectively, and the joint density of the event [T={, X(t)=a(t) | X(t 0 )=x 0 ] is g[a({), { | x 0 , t 0 ]_ f [a(t), t | a({), {]. Because this decomposition of the sample paths passing through (a(t), t) is exhaustive and mutually exclusive, the probability density associated with paths passing though this point may be obtained by summing across values of { to yield f [a(t), t | x 0 , t 0 ]= | t g[a({), { | x 0 , t 0 ] f [a(t), t | a({), {] d{. (35) t0 Further discussion of this equation and its ramifications may be found in Durbin (1971) and van Kampen (1992, pp. 307311). Equation (35) is formally a Volterra equation of the first kind, in which the unknown first passage time density g appears only under the integral sign. Although this equation is in principle soluble, the problem in attempting to work with it directly is that the kernel of the equation is (weakly) singular; that is, 10 When deriving (35) it is technically more correct to consider the probability density function, not of the point a(t), but of a(t)=lim = Ä 0 a(t)+=. Then { is always strictly less than t and the transition density f [a(t), t | a({), {] always well defined. But as the event {=t is of probability zero, the additional rigor of this refinement does little more than burden the notation and obscure the essential simplicity of the relationship expressed by (35). 432 PHILIP L. SMITH lim { Ä t f [x, t | y, {]=$(x& y), the Dirac delta function. The consequence of this divergence is, as discussed by Durbin (1971), that any numerical method that attempts to approximate the integral in (35) directly, as the limit of sums, will be inherently unstable. To circumvent this problem, recent research has sought to find ways to transform the equation to remove the singularity and thus allow the development of stable approximation methods. Most notably, Ricciardi and coworkers (e.g., Buonocore, Nobile, 6 Ricciardi, 1987; Buonocore, Giorno, Nobile, 6 Ricciardi, 1990; Giorno, Nobile, Ricciardi, 6 Sato, 1989; Ricciardi, Sacerdote, 6 Sato, 1983, 1984) have considered a variety of ways for stably transforming this equation that lead to practical computational methods. The most tractable and efficient of these methods, and hence the one that is most likely to be of value to cognitive modelers, was derived by Buonocore et al. (1987, 1990). This approach forms the basis of the methods described here. We follow Buonocore et al. (1987) in showing how to transform (35) into a Volterra integral equation of the second kind, in which the kernel goes to zero with t&{. We then describe a discrete time analogue of the resulting equation which lends itself readily to the development of an efficient numerical algorithm. Following Buonocore et al. (1987), we present these results in the form of a lemma followed by a theorem. In their work, they derived an expression for the kernel of the transformed equation that is somewhat more general than is needed for practical computational problems. Accordingly, we present a version of their results which is specialized for these applications. The reader is referred to the original article for further details. Lemma 1. Define .[a(t), t | y, {]= F[a(t), t | y, {], t (36) where F[a(t), t | y, {] is the transition distribution in (1) evaluated at the point a(t). Then g[a(t), t | x 0 , t 0 ] =&2.[a(t), t | x 0 , t 0 ]+2 | t g[a({), { | x 0 , t 0 ] .[a(t), t | a({), {] d{. (37) t0 Proof. Recalling the definition of the transition density in (26), we integrate over the state variable in the renewal equation (35) from a(t) to and then exchange the order of integration to obtain 1&F[a(t), t | x 0 , t 0 ]= = | | t g[a({), { | x 0 , t 0 ] t0 | f [x, t | a({), {] dx d{ a(t) t g[a({), { | x 0 , t 0 ][1&F[a(t), t | a({), {]] d{. t0 433 STOCHASTIC DYNAMIC MODELS Differentiating this equation with respect to t, making use of the elementary relation d dt | t h(t, {) d{=h(t, t)+ t0 h(t, {) d{ t0 t | t and the definition (36), yields &.[a(t), t | x 0 , t 0 ]= g[a(t), t | x 0 , t 0 ]&F[a(t), t | a(t), t] g[a(t), t | x 0 , t 0 ] | t g[a({), { | x 0 , t 0 ] .[a(t), t | a({), {] d{. & t0 By virtue of the relation 11 lim F[a(t), t | a({), {]= 12 , {Ät the second term on the right-hand side reduces to g[a(t), t | x 0 , t 0 ]2. The resulting equation may then be rearranged to yield g[a(t), t | x 0 , t 0 ]=&2.[a(t), t | x 0 , t 0 ] +2 | t g[a({), { | x 0 , t 0 ] .[a(t), t | a({), {] d{, t0 as asserted by the preceding lemma. Theorem 1. Define 9[a(t), t | y, {]=.[a(t), t | y, {]+k(t) f [a(t), t | y, {], (38) where k(t) is an arbitrary, continuous function of time, defined on [0, ), and .[a(t), t | y, {] is as defined in Lemma 1. Then g[a(t), t | x 0 , t 0 ]=&29[a(t), t | x 0 , t 0 ] +2 | t g[a({), { | x 0 , t 0 ] 9[a(t), t | a({), {] d{. (39) t0 11 This limit relation is attributed by Buonocore et al. (1987) to Fortet (1943). The property is rather unintuitive because it asserts that as the interval t&{ becomes small, P[X(t)>X({) | X({)=x] approaches P[X(t)<X({) | X({)=x] for all x in the interior of X, the state space of X(t), regardless of its infinitesimal moments. This property arises because in the SDE (6) and its various special cases, the deterministic change +(X(t), t) dt is of order dt, whereas the stochastic perturbation _(X(t), t) dB(t) is of order - dt. Because the stochastic part of dX(t) goes to zero more slowly than the deterministic part, the limiting transition distribution is determined by the (symmetrical) stochastic term only. 434 PHILIP L. SMITH Remark. Equations (37) and (39) are Volterra integral equations of the second kind, in which the unknown first passage time density g[a(t), t | x 0 , t 0 ] is defined jointly in terms of its values at preceding times, 0{<t, and the values of a known kernel function (.[a(t), t | y, {] or 9[a(t), t | y, {], respectively). The purpose of the greater generality of the kernel in (39) is that the arbitrariness of the function k(t) may be used to ensure that the kernel goes to zero with t&{. We derive expressions for the kernel function for specified diffusion processes subsequently. Proof. Substituting the definition of 9[a(t), t | y, {] from (38) into the righthand side of (39) yields g[a(t), t | x 0 , t 0 ]=&2.[a(t), t | x 0 , t 0 ] +2 | t g[a({), { | x 0 , t 0 ] .[a(t), t | a({), {] d{ t0 { +2k(t) & f [a(t), t | x 0 , t 0 ] +2 | t = g[a({), { | x 0 , t 0 ] f [a(t), t | a({), {] d{ . t0 The sum of the first and second terms on the right-hand side of this equation is g[a(t), t | x 0 , t 0 ], by Lemma 1. The integral in braces in the last line may be recognized as the right-hand side of the renewal equation (35) and is thus equal to f [a(t), t | x 0 , t 0 ], by virtue of that equation. The entire expression in braces in the second line of the preceding equation is therefore zero. Therefore the left- and righthand sides of the equation are equal, which completes the proof. With k(t) chosen in a manner to be determined subsequently, lim { Ä t 9[a(t), t | a({), {]=0. Under these circumstances, as discussed by Buonocore et al. (1987), the integral in (39) may be approximated numerically by a sum of the form g[a(t 0 +2), t 0 +2 | x 0 , t 0 ]=&29[a(t 0 +2), t 0 +2 | x 0 , t 0 ] g[a(t 0 +k2), t 0 +k2 | x 0 , t 0 ]=&29[a(t 0 +k2), t 0 +k2 | x 0 , t 0 ] k&1 +22 : g[a(t 0 + j2), t 0 + j2 | x 0 , t 0 ] j=1 _9[a(t 0 +k2), t 0 +k2 | a(t 0 + j2), t 0 + j2]. k=2, 3, ... (40) These expressions, which may be implemented easily on a personal computer, converge to the true first passage time density in (39) as the interval of approximation, STOCHASTIC DYNAMIC MODELS 435 FIG. 5. Renewal representation for the two-sided first passage time problem for a diffusion process X(t) through a pair of time-dependent absorbing barriers a 1(t) and a 2(t). Equations (41a) and (42b) give decompositions of f [a 1(t), t | x 0 , t 0 ] and f [a 2(t), t | x 0 , t 0 ], the probability densities associated with points on the upper and lower barriers, a 1(t) and a 2(t), at time t. The first barrier crossing is at a 1({) or a 2({) at time {, {t. The process then makes a transition to either a 1(t) or a 2(t) in time t&{. 2, becomes small. 12 The reader is referred to the original article for a convergence proof. The Two-Barrier Case The renewal approach described in the preceding section may also be applied to the solution of the first passage time problem for a diffusion process constrained by a pair of time-varying absorbing boundaries, as shown in Fig. 5. Buonocore et al. (1990) showed that the same approach to regularizing the kernel of the integral equation can be applied to the two-barrier problem to obtain a system of equations that provides the basis for a stable numerical approximation. Our presentation summarizes the main results from their article. To stress the similarity in the approaches to the one- and the two-barrier problems, we again present the results in the form of a lemma followed by a theorem. The basis for the solution of the two-barrier problem is a pair of simultaneous renewal equations, rather than the single equation considered previously. The probabilistic foundation for these equations is suggested by Fig. 5. As in the previous section, we let f (x, t | y, {) denote the free transition density of the diffusion process 12 For homogeneous Wiener and OU processes with time-varying boundaries, Buonocore et al. (1987) showed that convergence of (40) requires that the boundary a(t) be twice-differentiable, i.e., be of class C 2[0, ). The equivalent formulation for time-inhomogeneous processes with constant boundaries requires that m X (t; {), the mean of X(t) in (27) or (28), be of class C 2[0, ). In either instance, this degree of smoothness is required to obtain convergence of 9[a(t), t | a({), {] by L'Ho^pital's rule in the final step of Theorem 4. A model that fails to exhibit these smoothness requirements is the double halfwave rectifier model of Fig. 2a. There the functions m X (t; {) are only C 1[0, ) because the derivatives of [+(t)&c] + and [+(t)&c] & have one or more points of discontinuity. This problem can be circumvented by approximating these latter functions by functions with derivatives that are everywhere continuous. 436 PHILIP L. SMITH X(t) and let g 1[a 1(t), t | x 0 , t 0 ] and g 2[a 2(t), t | x 0 , t 0 ] denote the first passage time densities of X(t) through the absorbing barriers a 1(t) and a 2(t), respectively, where P[X(t 0 )=x 0 ]=1 and a 2(t 0 )<x 0 <a 1(t 0 ). We consider a decomposition of f [a 1(t), t | x 0 , t 0 ], the probability density of a point on the boundary a 1(t) at time t, given that the process started at the point x 0 <a 1(t 0 ) at time t 0 . Evidently, the process must have made at least one boundary crossing at some time {t. As shown in Fig. 5, this may occur in one of two ways. The process may first cross the upper boundary a 1({) at time { and then make a transition to the point a 1(t) in the interval t&{. The probability densities associated with the portions of the sample path before and after the first boundary crossing will then be g 1[a 1({), { | x 0 , t 0 ] and f [a 1(t), t | a 1({), {], respectively. Alternatively, the process may first cross the lower boundary a 2({) at time { and then subsequently make a transition to a point on the upper barrier a 1(t) at time t. The probability densities associated with the two portions of the sample path will then be g 2[a 2({), { | x 0 , t 0 ] and f [a 1(t), t | a 2({), {]. Because this decomposition according to time of first boundary crossing is exhaustive and mutually exclusive, and by virtue of the strong Markov character of X(t), the following relationship holds for sample paths crossing the upper boundary at time t: f [a 1(t), t | x 0 , t 0 ]= | t g 1[a 1({), { | x 0 , t 0 ] f [a 1(t), t | a 1({), {] d{ t0 | t g 2[a 2({), { | x 0 , t 0 ] f [a 1(t), t | a 2({), {] d{. + (41a) t0 An analogous decomposition of f [a 2(t), t | x 0 , t 0 ], the probability density of a point a 2(t) on the lower boundary at time t, gives a second equation f [a 2(t), t | a 0 , t 0 ]= | t g 2[a 2({), { | x 0 , t 0 ] f [a 2(t), t | a 2({), {] d{ t0 | t g 1[a 1({), { | x 0 , t 0 ] f [a 2(t), t | a 1({), {] d{. + (41b) t0 In principle, these equations may be solved simultaneously to yield the pair of unknown first passage time densities g 1[a 1(t), t | x 0 , t 0 ] and g 2[a 2(t), t | x 0 , t 0 ] but, as in the one barrier case, the singularity of f (x, t | y, {) as { approaches t precludes the development of a stable numerical approximation scheme based directly on these equations. Instead, a transformation is sought which removes the singularity. We have the following lemma. Lemma 2. Define .[a i (t), t | y, {], i=1, 2, as in Lemma 1. Then g 1[a 1(t), t | x 0 , t 0 ]=&2.[a 1(t), t | x 0 , t 0 ] +2 | t g 1[a 1({), { | x 0 , t 0 ] .[a 1(t), t | a 1({), {] d{ t0 +2 | t t0 g 2[a 2({), { | x 0 , t 0 ] .[a 1(t), t | a 2({), {] d{ (42a) STOCHASTIC DYNAMIC MODELS 437 and g 2[a 2(t), t | x 0 , t 0 ]=2.[a 2(t), t | x 0 , t 0 ] &2 &2 | | t g 1[a 1({), { | x 0 , t 0 ] .[a 2(t), t | a 1({), {] d{ t0 t g 2[a 2({), { | x 0 , t 0 ] .[a 2(t), t | a 2({), {] d{. (42b) t0 Proof. Integrating (41a) from a 1(t) to and (41b) from & to a 2(t) followed by an exchange in the order of integration gives 1&F[a 1(t), t | x 0 , t 0 ] = | t g 1[a 1({), { | x 0 , t 0 ][1&F[a 1(t), t | a 1({), {]] d{ t0 + | t g 2[a 2({), { | x 0 , t 0 ][1&F[a 1(t), t | a 2({), {]] d{ (43a) t0 and F[a 2(t), t | x 0 , t 0 ]= | t g 2[a 2({), { | x 0 , t 0 ] F[a 2(t), t | a 2({), {] d{ t0 | t g 1[a 1({), { | x 0 , t 0 ] F[a 2(t), t | a 1({), {] d{. + (43b) t0 We differentiate these equations with respect to t, making use of the limit relations lim F[a 1(t), t | a 1({), {]= 12 {Ät lim F[a 2(t), t | a 2({), {]= 12 {Ät lim F[a 1(t), t | a 2({), {]=1 {Ät lim F[a 2(t), t | a 1({), {]=0. {Ät The derivative of (43a) is &.[a 1(t), t | x 0 , t 0 ] = g 1[a 1(t), t | x 0 , t 0 ]& g 1[a 1(t), t | x 0 , t 0 ] F[a 1(t), t | a 1(t), t] & | t g 1[a 1({), { | x 0 , t 0 ] .[a 1(t), t | a 1({), {] d{ t0 + g 2[a 2(t), t | x 0 , t 0 ]& g 2[a 2(t), t | x 0 , t 0 ] F[a 1(t), t | a 2(t), t] & | t t0 g 2[a 2({), { | x 0 , t 0 ] .[a 1(t), t | a 2({), {] d{, (44) 438 PHILIP L. SMITH which, by virtue of the first and third equalities of (44), may be simplified and rearranged to give g 1[a 1(t), t | x 0 , t 0 ]=&2.[a 1(t), t | x 0 , t 0 ] +2 +2 | | t g 1[a 1({), { | x 0 , t 0 ] .[a 1(t), t | a 1({), {] d{ t0 t g 2[a 2({), { | x 0 , t 0 ] .[a 1(t), t | a 2({), {] d{, t0 which is (42a). Proceeding in a similar manner with (43b) yields .[a 2(t), t | x 0 , t 0 ] = g 2[a 2(t), t | x 0 , t 0 ] F[a 2(t), t | a 2(t), t] | t g 2[a 2({), { | x 0 , t 0 ] .[a 2(t), t | a 2({), {] d{ + t0 +g 1[a 1(t), t | x 0 , t 0 ] F[a 2(t), t | a 1(t), t] | t g 1[a 1({), { | x 0 , t 0 ] .[a 2(t), t | a 1({), {] d{, + t0 which, by virtue of the second and fourth equalities of (44) reduces to g 2[a 2(t), t | x 0 , t 0 ]=2.[a 2(t), t | x 0 , t 0 ] &2 &2 | | t g 1[a 1({), { | x 0 , t 0 ] .[a 2(t), t | a 1({), {] d{ t0 t g 2[a 2({), { | x 0 , t 0 ] .[a 2(t), t | a 2({), {] d{, t0 which is (42b), thus proving the lemma. Equations (42a) and (42b) are Volterra integral equations of the second kind, in which the unknown first passage time densities are defined in terms of their values at preceding times {<t and the kernel function .(x, t | y, {). We transform these equations to yield new integral equations whose kernels can be chosen in a way that they go to zero as { Ä t. Theorem 2. Set 9[a 1(t), t | y, {]=.[a 1(t), t | y, {]+k 1(t) f [a 1(t), t | y, {], (45a) 9[a 2(t), t | y, {]=.[a 2(t), t | y, {]+k 2(t) f [a 2(t), t | y, {]. (45b) and 439 STOCHASTIC DYNAMIC MODELS Then g 1[a 1(t), t | x 0 , t 0 ]=&29[a 1(t), t | x 0 , t 0 ] | +2 | +2 t g 1[a 1({), { | x 0 , t 0 ] 9[a 1(t), t | a 1({), {] d{ t0 t g 2[a 2({), { | x 0 , t 0 ] 9[a 1(t), t | a 2({), {] d{ (46a) t0 and g 2[a 2(t), t | x 0 , t 0 ]=29[a 2(t), t | x 0 , t 0 ] &2 | &2 t g 1[a 1({), { | x 0 , t 0 ] 9[a 2(t), t | a 1({), {] d{ t0 | t g 2[a 2({), { | x 0 , t 0 ] 9[a 2(t), t | a 2({), {] d{. (46b) t0 Proof. First consider (46a). Evaluating the right-hand side of this equation using the definition of 9[a 1(t), t | y, {] in (45a) yields g 1[a 1(t), t | x 0 , t 0 ]=&2[.[a 1(t), t | x 0 , t 0 ]+k 1(t) f [a 1(t), t | x 0 , t 0 ]] {| +2 t g 1[a 1({), { | x 0 , t 0 ] .[a 1(t), t | a 1({), {] d{ t0 | +k 1(t) {| +2 t g 1[a 1({), { | x 0 , t 0 ] f [a 1(t), t | a 1({), {] d{ t0 = t g 2[a 2({), { | x 0 , t 0 ] .[a 1(t), t | a 2({), {] d{ t0 +k 1(t) | t = g 2[a 2({), { | x 0 , t 0 ] f [a 1(t), t | a 2({), {] d{ . t0 The integrals with coefficients k 1(t) on the right-hand side of this equation may be recognized as the terms on the right-hand side of the renewal equation (41a). By virtue of the renewal equation, the sum of these terms, together with their coefficients, is k 1(t) f [a 1(t), t | x 0 , t 0 ]. With this substitution the preceding equation simplifies to g 1[a 1(t), t | x 0 , t 0 ]=&2.[a 1(t), t | x 0 , t 0 ] &2k 1(t) f [a 1(t), t | x 0 , t 0 ]+2k 1(t) f [a 1(t), t | x 0 , t 0 ] +2 + | {| t t g 1[a 1({), { | x 0 , t 0 ] .[a 1(t), t | a 1({), {] d{ t0 = g 2[a 2({), { | x 0 , t 0 ] .[a 1(t), t | a 2({), {] d{ . t0 440 PHILIP L. SMITH The sum of the second and third terms in this equation remaining terms is g 1[a 1(t), t | x 0 , t 0 ] by Lemma 2. The equals the right-hand side, which verifies (46a). Equation (46b) is proved analogously. The right-hand expanded using the definition of 9[a 2(t), t | y, {] in (45b) is zero; the sum of the left-hand side therefore side of the equation is to yield g 2[a 2(t), t | x 0 , t 0 ]=2[.[a 2(t), t | x 0 , t 0 ]+k 2(t) f [a 2(t), t | x 0 , t 0 ]] {| &2 t g 1[a 1({), { | x 0 , t 0 ] .[a 2(t), t | a 1({), {] d{ t0 | +k 2(t) {| &2 t g 1[a 1({), { | x 0 , t 0 ] f [a 2(t), t | a 1({), {] d{ t0 = t g 2[a 2({), { | x 0 , t 0 ] .[a 2(t), t | a 2({), {] d{ t0 +k 2(t) | t = g 2[a 2({), { | x 0 , t 0 ] f [a 2(t), t | a 2({), {] d{ . t0 The integrals with coefficients k 2(t) are the terms on the right-hand side of the renewal equation (41b). Summing these terms together with their coefficients gives k 2(t) f [a 2(t), t | x 0 , t 0 ], thereby reducing the equation to g 2[a 2(t), t | x 0 , t 0 ]=2.[a 2(t), t | x 0 , t 0 ] &2k 2(t) f [a 2(t), t | x 0 , t 0 ]+2k 2(t) f [a 2(t), t | x 0 , t 0 ] &2 & | {| t g 1[a 1({), { | x 0 , t 0 ] .[a 2(t), t | a 1({), {] d{ t0 t = g 2[a 2({), { | x 0 , t 0 ] .[a 2(t), t | a 2({), {] d{ . t0 In this equation, the second and third terms again sum to zero; the remaining terms sum to g 2[a 2(t), t | x 0 , t 0 ] by Lemma 2. The left- and right-hand sides of (46b) are thus equal, which proves the theorem. On the assumption that the kernel functions 9[a i (t), t | y, {] i=1, 2, may be chosen in such a way that they go to zero as { Ä t, the integrals in (46a) and (46b) may be approximated numerically by sums and the resulting equations solved simultaneously as follows: g 1[a 1(t 0 +2), t 0 +2 | x 0 , t 0 ] =&29[a 1(t 0 +2), t 0 +2 | x 0 , t 0 ] STOCHASTIC DYNAMIC MODELS 441 g 1[a 1(t 0 +k2), t 0 +k2 | x 0 , t 0 ] =&29[a 1(t 0 +k2), t 0 +k2 | x 0 , t 0 ] k&1 +2 2 : g 1[a 1(t 0 + j2), t 0 + j2 | x 0 , t 0 ] j=1 _9[a 1(t 0 +k2), t 0 +k2 | a 1(t 0 + j2), t 0 + j2] k&1 +2 2 : g 2[a 2(t 0 + j2), t 0 + j2 | x 0 , t 0 ] j=1 _9[a 1(t 0 +k2), t 0 +k2 | a 2(t 0 + j2), t 0 + j2]. k=2, 3... (47a) g 2[a 2(t 0 +2), t 0 +2 | x 0 , t 0 ] =29[a 2(t 0 +2), t 0 +2 | x 0 , t 0 ] g 2[a 2(t 0 +k2), t 0 +k2 | x 0 , t 0 ] =29[a 2(t 0 +k2), t 0 +k2 | x 0 , t 0 ] k&1 &22 : g 1[a 1(t 0 + j2), t 0 + j2 | x 0 , t 0 ] j=1 _9[a 2(t 0 +k2), t 0 +k2 | a 1(t 0 + j2), t 0 + j2] k&1 &22 : g 2[a 2(t 0 + j2), t 0 + j2 | x 0 , t 0 ] j=1 _9[a 2(t 0 +k2), t 0 +k2 | a 2(t 0 + j2), t 0 + j2]. k=2, 3... (47b) As in the single-barrier case, these equations may be evaluated straightforwardly on a personal computer. A numerical study of their convergence properties may be found in Buonocore et al. (1990). THE KERNEL OF THE INTEGRAL EQUATION In this section we see how to obtain the unknown functions k i (t) in Theorems 1 and 2 in such a way that the kernels 9[a i (t), t | a i ({), {] vanish as { Ä t. Together with the transition density f (x, t | y, {), these functions may be used in (41) and (47) to obtain the first passage time densities for X(t) numerically. We proceed to characterize the kernel function for those diffusions whose transition densities are Gaussian or which can be made Gaussian by appropriate transformation. This class of functions includes the general Gaussian process (14) and its various special cases, as well as the lognormal process (32). To this end, we seek to characterize those processes that can be transformed into a standard Brownian motion (Wiener) process by appropriate transformation of the time and state coordinates. A constructive proof of the conditions for the existence of such a transformation was provided by Cherkasov (1957) and recast into a convenient form for 442 PHILIP L. SMITH applications by Ricciardi (1976). A slightly more succinct statement of this result may be found in Ricciardi and Sato (1983). We record their result in the form of a theorem. Theorem 3 (Ricciardi 6 Sato, 1983). Let X(t) be a diffusion process satisfying the SDE (6) with drift +(x, t) and diffusion coefficient _ 2(x, t). Let _ 2$ x (x, t)= (x) _ 2(x, t) and _ 2$ t (x, t)=(t) _ 2(x, t) be the first partial derivatives of the diffusion coefficient with respect to its state and time coordinates, respectively. If there exists a pair of functions c 1(t) and c 2(t) such that +(x, t)= _ 2$ x (x, t) _(x, t) + c 1(t)+ 4 2 { | x c 2(t) _ 2( y, t)+_ 2$ t( y, t) dy , _ 3( y, t) = (48) then there exists a coordinate transformation, X(t) Ä X*(t*), of the form x*=9(x, t) (49) t*=8(t), such that X*(t*)=B(t*) is a standard Brownian motion. This transformation is x*=9(x, t) _ =exp & 1 2 | t c 2(s) ds &| x 1 dy & _( y, t) 2 | t _ c 1(s) exp & 1 2 | s & c 2(z) dz ds (50) t*=8(t) = | t _ | exp & t0 s & c 2(z) dz ds. Remark. The crucial requirement for the existence of this transformation is the existence of a pair of functions c 1(t) and c 2(t) of the time coordinate alone that relate the drift and diffusion coefficients in the manner indicated. Under this transformation, the new state coordinate x* is a function jointly of the old state coordinate and the old time coordinate; the new time coordinate is a function of the old time coordinate alone. Under the transformation X(t) Ä B(t*), the transition density of the process will be given by (27) with m X (t; {)=0 and _ 2 =1. Probabilistically, Theorem 3 asserts that the identity P[X(t)x | X({)= y]=P[B(t*)x* | B({*)= y*] holds with an appropriate choice of coordinates. By the previous remark P[B(t*)x* | B({*)= y*] x* = f B (x*, t* | y*, {*)= 1 - 2?(t*&{*) _ exp & (x*& y*) 2 , 2(t*&{*) & 443 STOCHASTIC DYNAMIC MODELS where a subscript on the transition density has been introduced to identify the process. Let 9 $x (x, t) and 9 $t(x, t) denote the first partial derivatives of 9(x, t) with respect to its first and second arguments, respectively. From the preceding equations, f X (x, t | y, {)= = P[X(t)x | X({)= y] x dx* P[B(t*)x* | B({*)= y*] x* dx = f B[9(x, t), 8(t) | 9( y, {), 8({)] 9 $x (x, t) = 1 - 2?[8(t)&8({)] { exp & [9 (x, t)&9 ( y, {)] 2 9 $x (x, t). 2[8(t)&8({)] = (51) Equation (51) provides an expression, in (x, t) coordinates, for the transition density of a process X(t) that satisfies the relation (48). By virtue of (51), one may work in whichever of the (x, t) or (x*, t*) coordinate sets is most convenient and express the results in the other coordinates by transformation. Buonocore et al. (1987) worked in the (x*, t*) coordinate set and provided explicit expressions for the kernel of the time-homogeneous Brownian motion and OU processes only. This was subsequently extended to a wider class of time-homogeneous processes by Giorno, Nobile, Ricciardi, and Sato (1989). These methods were used by Heath (1992) and Smith (1995) to remove time inhomogeneities in the drifts of Brownian motion and OU processes, respectively. 13 Although this method is often easy to implement computationally, one of its limitations, as pointed out by Gutierrez Jaimez, Roman Roman, and Torres Ruiz (1995), is that the image of the absorbing boundary a(t) in transformed coordinates is a*(t*)=9(a(t), t)=9[a[8 &1(t*)], 8 &1(t*)]. To obtain an explicit expression for the absorbing barrier in the new coordinate space requires that the function 8(t) that maps the old time coordinate to the new time coordinate be invertible. As Gutierrez Jaimez et al. (1995) noted, there are cases of interest in applications where this property does not apply. In these cases, 8(t) must be inverted numerically and the application of the method becomes cumbersome. Accordingly, in the remainder of this section we use the methods of Gutierrez Jaimez et al. (1995) and work in the original coordinate set. The following theorem and associated preamble are adapted from results proved in their article. 13 In Smith (1995) a heuristic argument was used to justify this transformation. It can be made rigorous by using the Ito^ calculus to show that the SDE which is satisfied by the transformed process is the defining equation for a zero-drift OU process. 444 PHILIP L. SMITH We use the definition (36) to obtain an expression for .[a(t), t | y, {] and, thereby, an expression for the kernel function 9[a(t), t | y, {]. By (51), the transition distribution of X(t) in (x, t) coordinates is F(x, t | y, {)= 1 | - 2?[8(t)&8({)] 9(x, t) { exp & & [s&9( y, {)] 2 ds, 2[8(t)&8({)] = which, by an obvious change of variables, becomes F(x, t | y, {)= 1 - 2? | z(x, t) exp(&` 22) d`, & where z(x, t)= 9(x, t)&9( y, {) - 8(t)&8({) . By definition (36), .[a(t), t | y, {]= d 1 dt - 2? | z(a(t), t) exp(&` 22) d`, & where the expression on the right is the total derivative of the integral, considered as a function of t. This evaluates to .[a(t), t | y, {] = 1 - 2? { exp & [9 (a(t), t)&9 ( y, {)] 2 2[8(t)&8({)] = 2[8(t)&8({)][9 $x (a(t), t) a$(t)+9 $t(a(t), t)]&[9 (a(t), t)&9 ( y, {)] 8$(t) _ . 2[8(t)&8({)] 32 Multiplying the right-hand side of this expression by 9 $x (a(t), t)9 $x (a(t), t) and making use of the definition of f X (x, t | y, {) in (51) gives .[a(t), t | y, {]= f X [a(t), t | y, {] { _ a$(t)+ 9 $t(a(t), t) [9 (a(t), t)&9 ( y, {)] 8$(t) . & 9 $x (a(t), t) 2[8(t)&8({)] 9 $x (a(t), t) (52) = As the consequence of this equation, we have the following easy theorem. Theorem 4 (Adapted from Gutierrez Jaimez et al. (1995)). Let .[a(t), t | y, {] be as defined in (36). Then with k(t)=& 1 9 $t (a(t), t) a$(t)+ , 2 9 $x (a(t), t) _ & (53) 445 STOCHASTIC DYNAMIC MODELS the limit relation lim 9[a(t), t | a({), {]= lim [.[a(t), t | a({), {]+k(t) f X [a(t), t | a({), {]]=0 {Ät (54) {Ät holds. Remark. This theorem provides the conditions required for the vanishing of 9[a(t), t | a({), {], the kernel of the integral equation in Theorems 1 and 2, as { Ä t. The condition (53) is that k(t) be defined as a function of the absorbing barrier(s) and of the ratio of the first partial derivatives with respect to time and with respect to state of the coordinate transformation that maps the particular diffusion process to a standard Brownian motion. The conditions for the existence of this latter transformation are given by Theorem 3. Proof. Let { h(t, {)= a$(t)+ 9 $t (a(t), t) [9(a(t), t)&9(a({), {)] 8$(t) & +k(t) . 9 $x (a(t), t) 2[8(t)&8({)] 9 $x (a(t), t) = As f X [a(t), t | a({), {] is singular at {=t, the vanishing of the kernel 9[a(t), t | a({), {] in (54) requires that h(t, t)=lim h(t, {)=0. {Ät With k(t) defined by (53) this implies { lim a$(t)+ {Ät 9 $t (a(t), t) [9(a(t), t)&9(a({), {)] 8$(t) =0. & 9 $x (a(t), t) [8(t)&8({)] 9 $x (a(t), t) = Equivalently, multiplying both sides of this equation by 9 $x (a(t), t)8$(t), noting from the definitions in (50) that for finite t this quantity is never singular and never zero, yields lim {Ät { 9 $x (a(t), t) a$(t)+9 $t (a(t), t) [9(a(t), t)&9(a({), {)] =0. & 8$(t) [8(t)&8({)] = (55) But t&{ [9(a(t), t)&9(a({), {)] 9(a(t), t)&9(a({), {) = , [8(t)&8({)] t&{ 8(t)&8({) which in the limit { Ä t is (ddt)(9(a(t), t))8$(t), which is the first term on the left-hand side of (55). The function h(t, {) is therefore zero at {=t, which means 446 PHILIP L. SMITH that at this point the kernel 9[a(t), t | a(t), t]=h(t, t) f X [a(t), t | a(t), t] is an indeterminate form 0 } . Repeated application of L'Ho^pital's rule shows this to be zero, thus proving the theorem. Corollary. As 2 Ä 0, the numerical integral equations (40), (47a), and (47b) with kernel function(s) 9[a(t), t | y, {]= f X [a(t), t | y, {] 2 { _ a$(t)+ 9 $t (a(t), t) [9(a(t), t)&9( y, {)] 8$(t) & 9 $x (a(t), t) [8(t)&8({)] 9 $x (a(t), t) = (56) converge to the first passage time densities g T [a(t), t | x 0 , t 0 ], g 1[a 1(t), t | x 0 , t 0 ] and g 2[a 2(t), t | x 0 , t 0 ]. The kernels for (47a) and (47b) are obtained from (56) with a(t) set equal to a 1(t) and a 2(t), in turn. Proof. This follows immediately on substituting (52) and (53) into the definition (38). It provides an expression for the kernel that is convenient for applications. EXAMPLES The time-inhomogeneous Brownian motion (Wiener) process. We consider the kernel of the integral equation for the diffusion process that satisfies the SDE (12), with drift +(x, t)=+(t), diffusion coefficient _ 2(x, t)=_ 2, and t 0 =0, x 0 =0. With these values of drift and diffusion coefficients, (48) is satisfied with c 1(t)=2+(t)_ and c 2(t)=0. Substituting these values in (50) shows that the transformation that maps the time-inhomogeneous Brownian motion to a standard Brownian motion is 9(x, t)= 1 x& _ _ | t +(s) ds & and 8(t)=t. In this example, the mapping of the time coordinate is the identity; that is, the transformation is of the state coordinate only. We have +(t) 9 $t (a(t), t)=& ; _ 1 9 $x (a(t), t)= ; _ 8$(t)=1. Using these functions in (56) then yields 9[a(t), t | y, {]= f [a(t), t | y, {] a(t)& y& t{ +(s) ds a$(t)&+(t)& , 2 t&{ { = (57) 447 STOCHASTIC DYNAMIC MODELS with f [x, t | y, {] given by (27). A simple calculation shows that lim { Ä t 9[a(t), t | a({), {]=0, as required. Indeed, when a(t)=a (constant) and +(t)=+, 9[a(t), t | a({), {]=0 for all t, {, {t. Under these circumstances, the kernel of the integral equation vanishes uniformly, so the integral on the right-hand side of (39) is zero for all t. The first passage time density for a time-homogeneous Brownian motion process through a constant absorbing barrier is therefore g(a, t | 0, 0)=&29(a, t | 0, 0) = a f (a, t | 0, 0). t \+ We have thus recovered a well-known result. First passage times for a homogeneous Brownian motion process through a constant boundary follow a Wald (or inverse Gaussian) distribution (e.g., Karlin 6 Taylor, 1975, p. 363), whose density function is related to the transition density of the unrestricted process in the manner shown. An analogous result was derived by Buonocore et al. (1987) for a zero-drift Brownian motion through a linear boundary. By an appropriate change of measure on the underlying process, effected via an application of the Girsanov theorem (Karatzas 6 Shreve, 1991, pp. 196197), these results may be seen to be equivalent. The time-inhomogeneous OU process. We obtain the kernel of the integral equation for the diffusion process that satisfies the SDE (13), with drift +(x, t)= +(t)&#x and diffusion coefficient _ 2(x, t)=_ 2. Equation (48) in Theorem 3 is satisfied with c 1(t)=2+(t)_ and c 2(t)=&2#. By (50), for t 0 =0, x 0 =0, the transformation that maps the process to a standard Brownian motion is 9(x, t)= 8(t)= e #tx 1 & _ _ | t e #s+(s) ds 1 [e 2#t &1]. 2# Unlike the previous example, for the OU process, both the time and the state coordinates change under the indicated mapping. For this choice of functions we have 9 $t (a(t), t)= e #t[#a(t)&+(t)] ; _ 9 $x (a(t), t)= e #t ; _ 8$(t)=e 2#t. Equation (56) then yields 9[a(t), t | y, {]= f [a(t), t | y, {] [a$(t)+#a(t)&+(t) 2 & 2# a(t)&e &#(t&{)y& 1&exp[&2#(t&{)] _ | t { e &#(t&s)+(s) ds, &= , (58) 448 PHILIP L. SMITH with f (x, t | y, {) given by (28). For the time-homogeneous case +(t)=+ (constant), this reduces to 9[a(t), t | y, {]= f [a(t), t | y, {] a$(t)+#a(t)&+ 2 { & 2 exp[&#(t&{)] [exp[#(t&{)](#a(t)&+)&(# y&+)] . 1&exp[&2#(t&{)] = (59) This expression for the kernel of the integral equation was derived by Buonocore et al. (1987) using a slightly different method. As in the preceding example, an easy calculation shows that lim { Ä t 9[a(t), t | a({), {]=0. Buonocore et al. also showed that the kernel vanishes uniformly for a hyperbolic absorbing boundary, in which case a simple, closed-form expression for the one-sided first passage time density can be found in a manner similar to that obtained for the Brownian motion process in the previous example. The interested reader is referred to their article for details. It should be noted that in the most common applications of the preceding equations the absorbing boundary (or boundaries) will be constant, and a$(t)=0. The lognormal process. As a final example of the use of these methods we obtain the kernel for a diffusion process satisfying the SDE (30) with +(x, t)=+(t) x, _ 2(x, t)=_ 2(t) x 2, t 0 =0, and X(0) # R + =(0, ). With these coefficients, the conditions of Theorem 3 are satisfied with c 1(t)=2+(t)_(t) and c 2(t)=&2_$(t)_(t). Equation (50) shows that the transformation that maps this process to a standard Brownian motion is 9(x, t)=log x& 8(t)= | t | t +(s) ds+ 12 | t _ 2(s) ds _ 2(s) ds. 0 For these expressions we have 9 $t (a(t), t)=&+(t)+ _ 2(t) ; 2 9 $x (a(t), t)= 1 ; a(t) 8$(t)=_ 2(t). Substitution of these functions in (56) yields a kernel of the form 9[a(t), t | y, {]= f [a(t), t | y, {] a$(t)&+(t) a(t) 2 { &_ 2(t) a(t) _ log a(t)&log y& t{ +(s) ds t{ _ 2(s) ds &= , (60) 449 STOCHASTIC DYNAMIC MODELS with f (x, t | y, t) given by (34). For the special case of _(t)=_ (constant), we recover the simpler expression derived by Gutierrez Jaimez et al. (1995), namely, 9[a(t), t | y, {]= f [a(t), t | y, {] a$(t)&+(t) a(t) 2 { &a(t) _ log a(t)&log y& t{ +(s) ds t&{ &= . (61) A calculation similar to that made in the two previous examples shows that for y=a({) these functions vanish as { approaches t, as required. MULTIVARIATE EXTENSIONS Application of the techniques described in the previous section will yield most of the statistics that are of interest in cognitive models. First passage time densities obtained in this way may be integrated numerically using the trapezoidal method or Simpson's rule (e.g., Dahlquist 6 Bjorck, 1974) to obtain values of the first passage time distributions G T (t), G 1(t), and G 2(t). These distributions may be evaluated at large values of t to estimate the absorption probabilities P[T<], for the single-barrier case, and P[T 1 <T 2 ], for the two-barrier case. When multiple processes X i (t), i=1, 2, ..., n, are involved, as occurs in the model of Fig. 2, the first passage time statistics of the resulting multivariate process are easily evaluated if the SDEs that define the constituent processes are uncoupled and the driving Brownian motion processes B i (t) are independent. Multivariate processes of this kind commonly arise in independent, parallel race models, in which the random variable of greatest interest is usually T=min[T 1 , T 2 , ..., T n ], (62) the time of the first-finishing process. If the random variables T 1 , ..., T n are each defined by an equation of the form (2) or (3), then we have the following geometrical interpretation of the random variable T: Let D1 /R n be defined by D1 =[x i : b i <x i <a i ; i=1, ..., n], where &b i <a i < and assume that X(0) # D1 . When the a i and b i are all finite, D1 is the interior of an n-dimensional hypercube. The random variable T then describes the time at which the vector-valued process X T (t)=[X 1(t), ..., X n(t)] first exits from the region D1 . (The superscript T in this notation represents matrix transposition.) Psychologically, this event corresponds to the time at which one of the coordinate accumulation processes X i (t) first exceeds its associated criterion or criteria. If g i (t) denotes the marginal first passage time density for the ith process, with g (n) T the first passage time density of T in (62), then we have the familiar expression n n g (n) T (t)= : g i (t) ` [1&G j (t)], i=1 j=1 j{i (63) 450 PHILIP L. SMITH (e.g., Ratcliff, 1978). Expressions of this form may be evaluated straightforwardly using the methods of the preceding section. When the processes X i (t) are coupled, as occurs in models with correlated noise, the transition distribution of the free process, X(t), may still be ascertained easily, but solution of the first passage time problem is appreciably more difficult. Under these circumstances, the evolution of X(t) is described by the vector-valued counterpart of (6), d X(t)=C(X(t), t) dt+D(X(t), t) d B(t), (64) where X(t) and B(t) are n- and m-dimensional random processes, respectively. The drift term, C(X(t), t), is an n-dimensional vector that is jointly a function of n state variables X i (t) and time t. The diffusion term, D(X(t), t), is an n_m matrix, each column of which is a function of the set of state variables and time. That m need not equal n in this equation is because the dimensionality of the process and the dimensionality of the superposed noise perturbations in general may be different. As was the case with the pair of scalar-valued SDEs (6) and (11), the most important special case of (64) is the one in which C(X(t), t) is linear in X(t) and D(X(t), t) is independent of state: d X(t)=[+(t)+A(t) X(t)] dt+_(t) d B(t). (65) Here +(t), A(t), and _(t) are (n_1), (n_n), and (n_m) matrix-valued functions, respectively. The function A(t) in this equation is a time-dependent coupling matrix, which represents statistical dependencies between the elements of X(t); +(t) is a stimulus-dependent forcing function (which in general is time-inhomogeneous), and _(t) is a time-dependent dispersion matrix, which represents statistical dependencies among components of the noise process. As in the case of the scalar-valued equation (11), the solution of (65) is relatively straightforward. Let the (n_n) matrix function 4(t) be the fundamental solution of the homogeneous, deterministic, first-order linear differential system 4$(t)=A(t) 4(t); 4(0)=I. As described in Hirsch and Smale (1974), a routine procedure for obtaining solutions to systems of equations of this kind is to find a coordinate transformation that uncouples the equations so that the matrix of the linear operator represented by A(t) is diagonal. The transformed system of equations can then be solved on an element-by-element basis and the results reexpressed in the old coordinates by inversion of the diagonalizing transformation. Once 4(t) has been obtained in this way, the solution to (65) may be expressed in terms of 4(t) and its inverse 4 &1(t) as follows (Karatzas 6 Shreve, 1991, pp. 354355), X(t)=4(t) | t 0 4 &1({) +({) d{+4(t) | t 0 4 &1({) _({) d B({), (66) 451 STOCHASTIC DYNAMIC MODELS where as usual we have assumed that P[X(0)=0]=1. That (66) solves (65) may be shown by an application of the multivariate form of the Ito^ transformation rule (e.g., Karlin 6 Taylor, 1981, p. 372). Evidently, the minimum requirement to ensure that the process X(t) in (66) is nondegenerate is that the matrix function 4(t) be invertible for all t. A more detailed characterization of how the properties of X(t) depend on A(t) and _(t) and a characterization of the conditions under which X(t) possesses a unique, stationary distribution may be found in Karatzas and Shreve (1991). In general, the process X(t) will have a multivariate normal distribution with mean m(t; {)=4(t) | t 4 &1(s) +(s) ds (67) { and time-dependent variancecovariance function v(t; {)=4(t) {| t = 4 &1(s) _(s) _ T (s)[4 &1(s)] T ds 4 T (t). { (68) In the special case in which A(t)=A and _(t)=_ the latter expression simplifies to t v(t; {)=exp(At) _| exp(&As) __ { T & exp(&A Ts) ds exp(A Tt). (69) As described by Smith (1996), the preceding equations provide a dynamic, multivariate generalization of SDT, or equivalently, of the General Recognition Theory of Ashby and Townsend (1986). These equations may be used to model performance in situations in which the observer's decision is based on a sampling interval of fixed duration. When the decision time depends on the observer sampling to a criterion, the decision time will be a random variable whose distribution is obtained by solving a first passage time problem for X(t). In principle, a renewal equation representation similar to (35) or (41) may be established (van Kampen, 1992, p. 311). The basis of this equation is suggested by Fig. 6. For concreteness, we assume that X T (t)=[X 1(t), X 2(t)] is a bivariate Gaussian processalthough the derivation we give applies in higher dimensions also. Let the transition density of X(t) be f (x, t | y, {). The precise form of this function will be determined by the mean and covariance functions in (67) and (68) with A(t), +(t), and _(t) appropriately specified. Let a(s) be a closed curve in R 2 parameterized by arc length and let D1 be the interior region of a(s), with X(0)=x 0 # D1 . Let D2 =R 2 &D1 be the complement of this region, where D2 includes the points on the boundary a(s). We define g[a(s), t | x 0 , t 0 ] to be the first passage time density of the process X(t) through the boundary a(s), where we stipulate as usual that P[X(0)=x 0 ]=1. Let x be a point on the boundary a(s). Then by an argument analogous to that used in the derivation of (35) we have f [x, t | x 0 , t 0 ]= | a(s) t g[a(s), { | x 0 , t 0 ] f [x, t | a(s), {] d{ ds. t0 (70) 452 PHILIP L. SMITH FIG. 6. Escape problem for a bivariate diffusion process X(t) from a closed region a(t). A general renewal representation of the first passage time density is given by Eq. (69). The outer integral in this equation is a contour integral evaluated around the boundary of the region D1 . Despite the similarities in the renewal equation representations for the scalar and vector cases, an important difference between (35) or (41) and (70) should be noted. Equations (35) and (41) involved, respectively, a single equation and a system of two equations, whereas (70) involves a system whose dimension is infinite. Such a system may be approximated by a system of finite dimension by approximating the boundary of D1 by a series of l linear segments 2a(i), i=1, 2, ..., l, and replacing the contour integral in (70) with a sum. A system of l simultaneous equations is then obtained with x i # 2a(i), which is in principle soluble for the first passage time density g[a(s), t | x 0 , t 0 ]. However, the high dimensionality of this system will make any numerical procedure that is based upon it expensive computationally. Under conditions identified by di Crescenzo, Ricciardi, Giorno, and Nobile (1991), first passage time problems involving multidimensional diffusion processes can sometimes be reduced to problems in a single dimension. This method can be applied to first passage time problems involving escape from a region bounded by an open (n&1)-dimensional, time-dependent surface in R n or by a pair of such surfaces. Under appropriate conditions, first passage problems of this form can be reduced to problems that involve the passage time for a one-dimensional process through a curve, a~(t), or a pair of such curves, a~ 1(t) and a~ 2(t). One model that can be dealt with using these methods is a stochastic, dynamic generalization of the integration model of Kinchla (1969, 1974), which has been investigated subsequently by various authors (e.g., Shaw, 1982; Smith, 1998a; see also MacMillan 6 Creelman, 1991). This model, which describes the detection of redundant signals in STOCHASTIC DYNAMIC MODELS 453 multidimensional displays, assumes that the observer monitors a set of n signal sources, the statistical characteristics of which are described by a set of Gaussian strength variables, X i , whose means depend on whether or not a source contained a signal. The integration model assumes that the observer sums these variables to create a composite decision variable i X i , which is compared to a criterion a to determine whether a detection response is made. Formally, _ n & P(S)=P : X i a , i where P(S) denotes the probability of a detection or ``Signal'' response. A dynamic generalization of this model can be obtained by assuming that the X i are coordinate processes of a multidimensional diffusion process and that the observer samples are from the display for t s time units. The observer responds ``Signal'' if and only if the sum of the coordinate processes X i (t) exceeds the criterion during this interval: _ { n = & P(S)=P inf t: : X i (t)a t s . i A two-boundary version of this model may be obtained by assuming that the sum of the coordinate processes is compared to a referent, c(t), and that the observer responds ``Signal'' or `` Noise'' depending on which of the boundaries a 1 or a 2 is first exceeded by the process i X i (t)&c(t). Here we state without proof the main results of di Crescenzo et al. (1991). For simplicity, we give results for a process in R 2 only. Proofs and a generalization to processes in R n may be found in the original article. We seek to identify conditions under which the first passage time density of the scalar-valued transformed-process Z(t)=![X(t)] through an absorbing barrier a~(t) is the same as that of X(t) through a curve a(t) in R 2, where a(t) is defined as follows: Assume that !(x) is monotone in x 1 and invertible, with inverse '(z, x 2 ). That is, z=!(x 1 , x 2 ) implies x 1 ='(z, x 2 ), and vice versa. Let a(t)=[(x 1 , x 2 ): x 1 = '(a~(t), x 2 )]. Let the bivariate diffusion process X(t) be defined by the SDE (64), with transition density f [x, t | y, {]= 2 P[X 1(t)x 1 , X 2(t)x 2 | X({)=y]. x 1 x 2 Let the transition density of the transformed process Z(t) be u(z, t | y, {). This density is obtained by integrating the transition density of X(t) over the region &<x 1 '(z, x 2 ), &<x 2 <, and then taking the derivative with respect to z to obtain: u(z, t | y, {)= | & f ['(z, x 2 ), x 2 , t | y, {] '(z, x 2 ) dx 2 . z (71) Now let g*['(a~(t), x 2 ), x 2 , t | x 0 , 0] dx 2 dt be the probability that X(t) first crosses the boundary a(t) in the interval (t, t+dt) through a linear element da(t) at the 454 PHILIP L. SMITH point ('(a~(t), x 2 ), x 2 ). The bivariate process X(t) will satisfy the following renewal equation f [x, t | x 0 , 0]= t & 0 | | g*['(a~({), x 2 ), x 2 , { | x 0 , 0] _ f [x, t | '(a~({), x 2 ), x 2 , {] d{ dx 2 . (72) This equation is a variant of (70), except that the curve a(t) is open instead of closed and, in addition, is permitted to depend on time. By definition, g[a(t), t | x 0 , 0]= | g*['(a~(t), x 2 ), x 2 , t | x 0 , 0] dx 2 . (73) & That is, the marginal first passage time density of X(t) through a(t) is obtained by integrating g* over all points in R 2 at which a boundary crossing may occur. Equation (72) may be transformed into a scalar-valued integral equation in z by integrating both sides of the equation over the region &<x 1 '(z, x 2 ), &<x 2 <, and then differentiating with respect to z, as indicated in (71). After exchanging the order of integration the resultant renewal equation is u[z, t | x 0 , 0]= t 0 & || g*['(a~({), x 2 ), x 2 , { | x 0 , 0] _u[z, t | '(a~({), x 2 ), x 2 , {] dx 2 d{. (74) If and only if u[z, t | y, {] has the representation u[z, t | y, {]=u~[z, t | `, {], (75) where `=!(y), then the term u[z, t | '(a~({), x 2 ), x 2 , {] may be taken out from under the inner integral in (74) by virtue of the fact that the set of points ('(a~({), x 2 ), x 2 ) falls on a locus of constant `, to give u[z, t | x 0 , 0]= | t g[a({), { | x 0 , 0] u~[z, t | a~({), {] d{, (76) 0 where the definition (73) has been used to eliminate the integral over x 2 in the final step. Equation (76) is a scalar-valued renewal equation which can be evaluated using the methods described previously (e.g., (40)). Analogous expressions may be developed for the first exit time for a region bounded by a pair of time-dependent curves a 1(t), a 2(t), each of which satisfies the conditions on a(t) described previously. In this case, the resultant renewal equation representation can be reduced to a pair of simultaneous scalar-valued equations that can be solved using the method of (47). The conditions required to reduce a multidimensional first passage time problem to a problem in one dimension are first, that the image of the curve a(t) under the 455 STOCHASTIC DYNAMIC MODELS mapping ! be invertible and second, that the transition density u has the representation (75), in which the form of the density depends only on the value of `=![( y 1 , y 2 )] and not on the values of y 1 and y 2 individually. The simplest and most tractable case in which these conditions are met is that of the first passage time for a Brownian motion process through a constant linear boundary. This model is a realization of the dynamic integration model described previously. We illustrate the condition (75) for the simplest case of an independent, unit-variance, bivariate Brownian motion with zero drift. Application to homogeneous (nonzero drift) Brownian motion and OU processes are described in di Crescenzo et al. (1991). In this case, X(t) has transition density f [x, t | y, {]= 1 2 (x i & y i ) 2 . exp & i=1 2?(t&{) 2(t&{) _ & The image of x under transformation is z=!(x)=x 1 +x 2 , so with a~(t)=a (constant), the absorbing boundary is the set of points a(t)/R 2, a(t)=[('(a, x 2 ), x 2 )]= [(a&x 2 , x 2 )]. Substituting for z in the previous equation, completing squares in the exponent, and using the result in (71) gives u(z, t | y, {)= 1 2?(t&{) 1 2x 2 &(z& y 1 + y 2 ) 2 - 2(t&{) _ { 1 z&( y + y ) & dx , 2 { - 2(t&{) = & | exp & & = 2 2 1 2 2 by virtue of the fact that '(z, x 2 )z=1. This expression may be recognized as the integral of a product of independent Gaussian densities in 2x 2 and z, respectively, each with standard deviation - 2(t&{). We may therefore integrate over x 2 to obtain the marginal density of z: u(z, t | y, {)= 1 z&( y 1 + y 2 ) exp & 2 2 - ?(t&{) - 2(t&{) 1 _ { 2 = &. Since `= y 1 + y 2 , we see that u(z, t | y, {) has a representation of the form u~(z, t | `, {), as required. Remark. Because of the requirement that the transformation z=!(x) be invertible, this method cannot be used to reduce the problem of escape from a closed region, as described in (70), to a problem in one dimension in this way. When certain strong symmetry conditions are met, the problem of escape from a closed region in R n can sometimes be reduced to a problem in one dimension by considering the Euclidian distance process Z(t)=- ni=1 X 2i (t). For the Brownian motion and OU processes, the Euclidian distance processes are known as the Bessel process and the radial OU process, respectively, both of which have well-defined transition densities on the positive real line R + =(0, ) (Karlin 6 Taylor, 1975, 456 PHILIP L. SMITH pp. 365371; 1981, pp. 333338). The first exit time for a zero-drift process from a closed, spherical region in R n is an example of a problem that can be analyzed in this way. Unfortunately, the processes of greatest interest in cognitive applications are those in which the drift term is nonzero (i.e, +(t){0), and in these cases, the symmetry conditions required to reduce the problem to a single dimension are absent. In principle, first passage time problems for diffusion processes may always be approximated using a discrete time, discrete state space representation, in which the problem is formulated as a first passage time problem for a finite state Markov chain (e.g., Karlin 6 Taylor, 1975; Bhattacharya 6 Waymire, 1990). First passage time distributions for such processes may be obtained even in the presence of time inhomogeneity, but their solution requires a matrix multiplication at each time step and is thus computionally expensive. Results similar to those described here for reducing a multidimensional first passage time problem to a problem in one dimension appear to have been obtained independently by Ashby and Schwartz (1996). Other relevant results are given in di Crescenzo, Giorno, Nobile, and Ricciardi (1995). SUMMARY Stochastic accumulation processes are important elements of many models in sensory and cognitive psychology, their role being to provide a theoretical foundation from which response time and accuracy predictions may be derived. These processes represent an essential link between observed performance, which is inherently probabilistic, and the underlying psychological mechanisms from which it arises. In the literature, the accumulation processes that have been proposed have been of two main kinds. The first, exemplified by SDT, assumes that the sampling time is determined by factors external to the information sample; the second, exemplified by the large and varied class of sequential sampling models, assumes that the sampling time is determined by the statistical properties of the sample itself. These two sorts of model lead, respectively, to the study of the free transition distribution of the unbounded accumulation process and to the study of its first passage time statistics. This article has investigated the foundations of a class of stochastic accumulation models that can be formulated as Markov processes in continuous time and continuous state space. The dynamics of information accrual in these models are represented by first-order, linear SDEs, whose solutions are diffusion processes. In general, the assumptions embodied in the defining SDEs of such models result in processes that are both temporally and spatially inhomogeneous. This article has provided a characterization of accumulation processes that either are Gaussian or can be made Gaussian by transformation of the time scale and state space. Methods were described for obtaining the time-dependent distribution of the unbounded accumulation process and for obtaining the first passage time distributions through either one or two absorbing barriers, which may themselves vary with time. For processes in one dimension, it was possible to give a fairly complete characterization both of the free transition distribution of the process and of its first passage STOCHASTIC DYNAMIC MODELS 457 time distribution. For processes in more than one dimension, a fairly complete characterization of the unbounded accumulation process was again possible, but the first passage time problems that arise in relation to processes of this kind are far less tractable. In important special cases, multidimensional first passage time problems can be reduced to equivalent problems in a single dimension. The conditions under which such a reduction is possible and the kinds of psychological processes that might be represented by models of this kind were described. APPENDIX We use the Ito^ calculus to obtain the solutions to the SDEs (10) and (30). The method of solution is adapted from Gardiner (1985, pp. 112113). Homogeneous case. We show that X(t)=X(0) U(t) solves (30) with X(0) the random initial value of X(t), X(0) # R + =(0, ), and U(t) defined as in (31). We write (30) in the form dX(t)=[b(t) dt+c(t) dB(t)] X(t), (A1) and consider the function Y(t)= f [X(t)]=log X(t). (A2) For functions of a single variable the Ito^ transformation formula (29) may be written in the form df (X(t))= f $(X(t)) dX(t)+ 12 f "(X)(dX(t)) 2. (A3) We recall that when working with stochastic differentials the following order relations hold: [dB(t)] 2 tdt; dB(t) dtt0; (dt) 2 t0. (A4) With f (x) defined as in (A2) we have f $(x)=1x and f "(x)=&1x 2. Applying the transformation formula (A3) to (A2) therefore yields dY(t)= dX(t) [dX(t)] 2 & . X(t) 2X 2(t) We substitute for dX(t) from (A1), noting that the order relations yield [dX(t)] 2 = c 2(t) X 2(t) dt, to obtain dY(t)=b(t) dt+c(t) dB(t)& 12 c 2(t) dt. 458 PHILIP L. SMITH This equation is integrable, i.e., Y(t)= | t b({) d{+ 0 | t 0 c({) dB({)& 12 | t c 2({) dt+Y(0), 0 which may be combined with the definition (A2) to give X(t)=X(0) exp _| t b({) d{+ 0 | t 0 c({) dB({)& 12 | t 0 & c 2({) dt , which is (32). Inhomogeneous case. We show that X(t) as defined in (33) solves (10), with U(t) defined as before. We seek a solution of the form X(t)=Z(t) U(t), (A5) where the function Z(t) is to be determined. Before proceeding, we use the Ito^ transformation rule to obtain the differential of the function U(t). To this end, we write this function as U(t)= f [V(t)]=exp V(t), where V(t) is the exponent in (31). The Ito^ rule (A3) applied to this function yields dU(t)=dV(t) e V(t) + 12 [dV(t)] 2 e V(t) =[b(t) dt+c(t) dB(t)] U(t). Also, the order relations (A4) yield [dU(t)] 2 =c 2(t) U 2(t) dt. Having established these preliminary facts, we write (A5) in the form Z(t)=X(t)[U(t)] &1. The chain rule form of the Ito^ formula (Karatzas 6 Shreve, 1991, pp. 150) applied to this function gives dZ(t)=dX(t)[U(t)] &1 +X(t) d[U(t)] &1 +dX(t) d[U(t)] &1. (The stochastic chain rule is derived using a procedure similar to that used to obtain (29): The product of two functions is expanded in a Taylor series, discarding terms of order (dt) 2 and above, and the result expressed in differential form.) The one-variable form of the Ito^ formula (A3) applied to the function d[U(t)] &1 gives d[U(t)] &1 =& dU(t) [dU(t)] 2 + . 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